ISSN 2520-5749

Proceedings
ICNS 2016
First International Conference of Natural Science
(ICNS 2016)
Organized By:
College of Basic Education, Charmo University, Chamchamal,
Sulaimani, Kurdistan-Iraq
11th − 12th July 2016
Royal Hotel, Qlyasan, Sulaimani, Kurdistan-Iraq
 ISSN 2520-5749

Proceedings
ICNS 2016
First International Conference of Natural Science
(ICNS 2016)

Organized By:
College of Basic Education, Charmo University, Chamchamal,
Sulaimani, Kurdistan-Iraq

11th − 12th July 2016

Royal Hotel, Qlyasan, Sulaimani, Kurdistan-Iraq
 ISSN 2520-5749

Proceedings
ICNS 2016
First International Conference of Natural Science
(ICNS 2016)

Editors

Ali W. Kareem Sangawi

Charmo University, Kurdistan

Salah Raza Saeed

Charmo University, Kurdistan-Iraq

Shwan Kamal Rachid

Charmo University, Kurdistan-Iraq
 Preface

ICNS2016 is the 1st International Conference of Natural Science which ran from 11th -12th July
and hosted by Charmo University. The conference is one of the most prestigious annual conferences
which offers more global scientific program and provides an excellent platform for national and
international prestige scholars in abroad to explore their various academic, latest research findings,
scientific activities, and innovations.
The call for paper resulted in a total of 200 submissions from around the world and almost
100 researchers from different Universities of Kurdistan Region, South part of Iraq and countries
including United State of America, Germany, Netherland, Iran, China, and France have attended the
ICNS2016. They were engaged with many scientists from across the world on these topics, and
shared the latest thinking on their researches and publications.
This conference had a multidisciplinary focus, addressing issues of Natural Science (Biology,
Chemistry, Physics, Mathematic, Geology, Computer Science, in addition Agriculture, and health
science which impact on the lives of Human beings every day.
The conference is started with plenary session, where keynote speeches were given by interna-
tional distinguished experts. Through the two days of the conference three parallel oral sessions
were held in separate halls. Posters are displayed throughout the two conference days. The cur-
rent proceeding comprises written contributions of contents of many selected presentations during
ICNS2016.
We gratefully acknowledge the ministry of higher education for the financial support, KAT
Group, and Kurd Oil for sponsoring the conference. We would like to thank all who contributed
to success the conference in particular, the organizing committee, and the program committee for
their competent evaluation of the large number of submissions, and their efforts in developing such
a stimulating and interesting conference program, without them we would not have been able to
organize this great meeting.

Sincerely,
Shwan K. Rachid
ICNS-2016 Chairman
Contents

ICNS 2016 Committees 3

Biology Session 5
Genetic Analysis in Five Inbred Line of Maize (Zea mays L.) and their F1 and F2 Hybrids
Under Different Conditions Using Line × Tester Analysis
Tawfiqa Sherwan Esmael, Hussaina Runak Ahmad, Hama-Ali Emad Omer∗ , Hamma-
Umin Bestoon Omer and Abdulkhaleqa Dana Azad . . . . . . . . . . . . . . . . . . 5
Advances in Detecting Breast Mass by 2D Ultrasound
Heamn Noori Abduljabbar and Sardar Yaba Perxdr . . . . . . . . . . . . . . . . . . . 27
Antimicrobial Activity of Silver Nanoparticles Synthesized From Banana’s Peel Extract
Huda Z. Majeed and Rana H. Al-Shammari . . . . . . . . . . . . . . . . . . . . . . . 35
The Effect of Cutting on Yield and Its Components of Six Genotypes of Barley Under
Slaimani Condition
Shara Jalal Hama, Shang Haseeb Abdulqader, Bestoon Omare and Kazal Rashed
Ahmad . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 43
Foliar Application of Humic Acid, Iron and Sprays Number on Chemical Quality of Grape
(Vitis vinifera L.) cv. TAIFI∗
Shawkat Mustafa Al-Atrushy1 and Shayan Akram Mustafa2 . . . . . . . . . . . . . 49
Induction of Sister Chromatid Exchange in Lymphocytes of Brick Kiln Workers in Bagh-
dad
Zainab M.T.Jaafar1 , Nazar A.Auda2 and Sa’adyia O. Mohmmed2 . . . . . . . . . . 59
Nutritional Status in Relation to Intestinal Infection Among Primary Schoolchildren in
Erbil Province Kurdistan-Iraq
Abdullah A. Hama1,2 and Zohair I.F. Rahemo3 . . . . . . . . . . . . . . . . . . . . 67

Chemistry Session 73
Evaluation of Hassira Crude Oil in Sarqalah, Kifri, Sulaimani, Kurdistan Region
Hersh J. Noori and Nabil A. Fakhre . . . . . . . . . . . . . . . . . . . . . . . . . . . 73
Batch and Flow Injection Spectrophotometric Determination of Amoxicillin in Pharma-
ceutical Formulations
Kamal M. Mahmoud* and Bery M. Rahman . . . . . . . . . . . . . . . . . . . . . . 81

Physics Session 99
Simulation of Cu Electrodeposition from Aqueous Solution for Solar Absorption Using
COMSOL Model
Azeez Abdullah Azeez Barzinjy, Haidar Jalal Ismae, Kadhim Qasim Jabbar, Mo-
hammed Abdullah Hamad and Samir Mustafa Hamad . . . . . . . . . . . . . . . . . 99
Neutronic and Thermal - Hydraulic Calculations for The Conversion of BN-350
Jamal Mohammed Rashid Abda and Abdul Karim Thamir Mohammed . . . . . . . . 109
Isolation, identification and determination of Amygdalin from the sweet and bitter almond
kernel fruit in Erbil city
Kamal M. Mahmoud* and Roza T. Yasin . . . . . . . . . . . . . . . . . . . . . . . . 115

1
 Density-Functional Calculations on Structural Properties of Sn Bulk and Nanocrystals
Botan Jawdat Abdullah, Musafa Saeed Omar and Qing Jiang . . . . . . . . . . . . . 127

Geoscience Session 133

Assessment and Validation of Landslide Susceptibility Mapping in Mountainous Basin
Using Spatial-Based Statistical Models and Remote Sensing Techniques
Himan Shahabi and Bakhtyar Ali Ahmad . . . . . . . . . . . . . . . . . . . . . . . . 133

Mathematics Session 147

αγ -Connectedness and Some Properties of α(γ,β) -Continuous Functions
Alias B. Khalaf and Hariwan Z. Ibrahim . . . . . . . . . . . . . . . . . . . . . . . . 147
Maximal λc -Open Sets
Halgwrd M. Darwesh, Sarhad F. Namiq and Wria K. Kadir . . . . . . . . . . . . . . 157
New R0 and R1 in Bitopological spaces
Hardi N. Aziz . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 163
Centralizers of Centrally Semiprime Rings
Adil Kadir Jabbar and Farhad Rafiq Krush . . . . . . . . . . . . . . . . . . . . . . . 171

Computer Session 181

Design Security System Based on AES and MD5 for Smart Card
Raghad Z.Yousif, Shahab.W.Kareem and Ammar. O.Hasan . . . . . . . . . . . . . . 181

2
 INTERNATIONAL ADVISORY COMMITTEE
Dr. Francis Burns
Dine College,
United States of America.
Dr. Stephan Hüttel
Fermentation Platform,
Helmholtz Centre for Infection Research,
Braunschweig, Germany.
Assist. Prof. Dr. Abdollah Hassanzadeh
Department of Physics,
University of Kurdistan, Sanandaj, Kurdistan, Iran
• ORGANIZING COMMITTEE
– Prof. Dr. Salah Raza Saeed
President of Charmo University
– Prof. Dr. Shwan Kamal Rachid
Vice President for Scientific affair of Charmo University
– Assist. Prof. Dr. Baram Ahmed
Vice President Office for Administration and finance of Charmo University
– Assist. Prof. Dr. Ali Wahab Kareem Sangawi
Dean of College of Basic Education of Charmo University
– Dr. Dana Khdr Sabir
• Scientific & Technical Committee
– Dr. Omed Mohammed Mustafa
– Dr. Ibrahim Sulaiman Hamad
– Dr. Rawand Bakir Noori Jaff
– Dr. Hemin Jalal Abdoul
– Dr. Shelanah Mohammed Raoof Salih
– Jwan Taher Rawuf
– Aso Mohammed Aladdin
– Zrar Khalid Abdul
– Beston Mohammed Qadir
– Ata Omer Salih
– Aziz Mohammad Abdulla
– Laila Ibrahim Faqe Saleh
– Mzhda Sabir Abdulkarim

3
4
Biology Session

Genetic Analysis in Five Inbred Line of Maize (Zea mays L.) and their F1 and
F2 Hybrids Under Different Conditions Using Line × Tester Analysis
Tawfiqa Sherwan Esmael, Hussaina Runak Ahmad, Hama-Ali Emad Omer∗ ,
Hamma-Umin Bestoon Omer and Abdulkhaleqa Dana Azad1
Field Crops Department, Faculty of Agricultural Sciences, University of
Sulaimani, 46001 Sulaimani, Kurdistan-Iraq
[email protected]
Abstract
Five inbred lines of maize (Zea mays L.) were used in this investigation to pro-
duce two generations. The hybridization consisted of 2 lines (MSI 4218 and MSI
4279) and 3 testers (MSI 42100, ZP 434 and 5012). Pooled analysis of variance
for both generations across locations confirmed the presence of highly significant
differences for all of the effects for most of the studied characters. The cross 2 × 5
gave the highest values for kernel yield/plant for F1 at Kanipanka and F2 at Qlyasan
locations and the cross 2 × 4 recorded the maximum values for kernel yield/plant
for F1 at Qlyasan and F2 at Kanipanka locations. The maximum positive GCA
effects were recorded by parent 2 and 3 for the first generation at Kanipanka and
Qlyasan locations respectively, while parent 4 obtained the highest values for the
second generation at both locations. Maximum positive SCA values obtained by
the cross 2 × 5 for F1 at Kanipanka and F2 at Qlyasan locations, while the cross
2 × 4 for F1 at Qlyasan and F2 at Kanipanka locations. The interaction of lines
× testers recorded the greatest contribution to the total variance concerning kernel
yield/plant for both generation and at both locations.
Keywords: Combining ability, Line × tester, Pooled analysis, Zea mays L.

Introduction

Maize (Zea mays L.) is widely cultivated crop throughout the word and has a
remarkable place among cereals and it is used as human food, animal feeding and
in industry [15]. The main objective of maize breeders is to obtain new inbred lines
and hybrids that will over dominance the existing hybrid due to kernel yield as the
most economically important characters in maize [25].
Line × tester analysis is one of breeding strategies to estimate combining ability
effects of genotypes and also to supply information concerning genetic mechanisms
controlling the characters [23]. Line × tester mating design provided dependable
information on the general and specific combining ability effects of parents and their
hybrid combinations in applied breeding programs [14]. The design has been widely
used in maize breeding by several workers and continues to be applied in quantitative
genetic studies in maize due to its significance [22]. Line × tester is useful in
deciding the relative ability of female and male lines to produce desirable hybrid
combinations. It also provides information on genetic components and enables

5
 the breeder to choose appropriate breeding methods for hybrid variety or cultivar
development programs. Information on combining ability effects helps the breeder
in choosing the parents with high general combining ability and hybrids with high
specific combining ability [8]. Line × tester was studied by [2; 4; 5; 10; 18; 19].
The aims of this study are to estimate the general combining ability for parents
of maize, the specific combining ability for hybrids of first and second generations
to select the parents with good GCA and crosses with good SCA effects through
line × tester analysis over different locations through evaluating three genotypes this
used as testers, and estimate the heterosis as percentage mean deviation of F1 ’s and
F2 ’s hybrid from mid parental values, also estimate proportional of contribution of
lines, testers and their interactions to total variance in the appeared characteristics.

Material and Methods

Five inbred lines of maize (Zea mays L.) namely (MSI 4218, MSI 4279, MSI 42100,
ZP 434 and 5012) were introduced from the Ministry of Agriculture, Sulaimani
Research Station, Bakrajo, Sulaimani, Iraq. The hybrids of the first generation were
generated during (autumn 2009) at Qlyasan location, the F1 ’s were sown and selfed
in the coming planting season (spring, 2010), to produce seed for F2 generation
at the same location. The hybridization consisted of 2 inbred lines namely (MSI
4218 and MSI 4279) which were used as female parents (lines), together with 3
inbred lines namely (MSI 42100, ZP 434 and 5012) which were used as male
parents (testers). The 6 cross hybrids for F1 ’s and F2 ’s with their 5 parents were
evaluated during the growing season (spring 2011) following line x tester analysis
at two environmentally different locations in Sulaimani, viz. Kanipanka Nursery
Station, Sulaimani Agricultural Directorate, Ministry of Agriculture (Lat 35o 22’
; N, Long 45o 43’ ; E, 550 masl) in Shahrazoor valley 35 Km east of Sulaimani
city and Qlyasan Agricultural Research Station, Faculty of Agricultural Sciences,
University of Sulaimani (Lat 35o 34’ 307” ; N, Long 45o 21’ 992” ; E, 765 masl),
2 Km north west of Sulaimani city. Irrigation, fertilization, and weed control were
accomplished according to normal field practices throughout the growing season.
Hills were overplanted and thinned after emergence for a final plant density of
about 55,000 plants / ha. Each cross was plants in four rows, 0.75 m apart and 5
m long with 0.25 m between plants. At full maturity, the central two rows of each
plot (excluding the boarders and 0.5 m from each end) were harvested manually
in both locations. Measurements were taken on the following parameters (average
of 5 plants/plot): plant height (cm), no. of ears/plant, no. of rows / ear, no. of
kernels / row, no. of kernels/ear, kernels weight/row (g), kernels weight/ear (g),
300-kernel weight (g), kernels yield/plant (g). Statistical analyses were carried out
using XLSTAT-Pro, Version 5.1. In both locations, the experiments were arranged
in a Completely Randomized Block Design (CRBD) with three replications, and
combined analyses of variance across locations were performed for both generations.
The following genetic parameters were computed according to [1; 23]:
1. Components of variance for both general and specific combining abilities,
2. General and specific combining ability effect,

6
 3. Standard error for the combining ability effects,
4. Heterosis,
5. Proportional contribution of lines, testers and their interaction to total variance.

Results and Discussions

Table 1 represents the pooled analysis of variance for combining ability for F1
generation at both locations. Environment effect was found to be highly significant
due to all of the studied characters except for the character no. of rows/ear which
was not significant. The mean squares due to genotypes and to the interaction of
genotypes × locations were highly significant except for no. of rows/ear which was
significant only for genotypes × locations interaction. ParentŠs mean square were
highly significant for all characters, and the mean square were highly significant for
the interaction of parents × locations with exception for the characters plant height
and no. of rows/ear which were found to be significant and no. of ears/plant which
was not significant. Crosses and the interaction of crosses × locations mean squares
were highly significant for all characters except no. of ears/plant for the crosses
which was significant only and no. of rows/ear for crosses and the interaction of
crosses × locations were not significant. The mean squares for the parents vs.
crosses and parents vs. crosses × locations were highly significant for all of the
studied characters, but the interaction of parents vs. crosses × locations was not
significant for the character kernel yield/plant. Lines mean squares were highly
significant for all characters except plant height and kernel weight/row which were
significant and no. of ears/plant and no. of rows/ear which were not significant.
TesterŠs mean squares were highly significant for all characters except the character
no. of rows/ear which was not significant, while for the component testers ×
locations were highly significant for all characters except no. of ears/plant which
was significant and no. of roes/ear and kernel weight/row were not significant. The
interactions of lines × testers and lines × testers × locations were highly significant
for all characters except the characters no. of ears/plant and no. of rows/ear which
were not significant for both components and the character kernel weight/row which
was significant only for lines × testers component.
The pooled analysis of variance for combining ability for F2 generation at both
locations represented in table 2. Environment mean squares were found to be
highly significant for the characters no. of kernels/row; kernel weight/ear and kernel
yield/plant only, while it was not significant for the rest. The mean squares due to
genotypes were highly significant for all characters except no. of ears/plant which
was not significant, while for the interaction of genotypes × locations were highly
significant also except the characters plant height and no. of ears/plant which were
not significant. ParentŠs mean squares were highly significant for all characters
except no. of ears/plant and kernel yield/plant which were not significant, but for
the interaction of parents × locations were highly significant for the characters no.
of rows/ear, no. of kernels/row, kernel weight/row and kernel yield/plant, but not
significant for plant height and no. of ears/plant. Regarding the mean squares
of crosses and the interaction of crosses × locations mean squares were highly

7
 significant for all characters except no. of ears/plant which was not significant.
The mean squares for the parents vs. crosses and parents vs. crosses × locations
were highly significant for all of the studied characters except no. of ears/plant
for parents vs. crosses which was significant and kernel yield/plant for parents
vs. crosses × locations which was significant also, and the characters plant height,
no. of ears/plant and kernel weight/row were not significant for the component
parents vs. crosses × locations. Lines mean squares were highly significant for all
characters except plant height, no. of ears/plant and no. of rows/ear which were not
significant, while for lines × locations were highly significant also for all characters
except plant height, no. of ears/plant and 300-kernel weight were not significant.
The mean squares due to testers and testers × locations were highly significant for
all characters except the character no. of ears/plant and kernel weight/row for testers
and plant height, no. of ears/plant and 300-kernel weight were not significant for
the component testers × locations. The mean squares calculated for lines × testers
were highly significant for all characters except no. of rows/ear which was significant
and no. of ears/plant which was not significant. The interactions of lines × testers
× locations were highly significant for all characters except no. of and no. of
rows/ear which was significant and plant height and no. of ears/plant which were
not significant
Data in table 3 explain the significance of the studied characters according to
their means squares values of Kanipanka and Qlyasan locations and both F1 and F2
generations. Respect to the mean squares due to genotypes at Kanipanka location
for the first generations it was highly significant for the characters plant height, no.
of ears/plant, no. of rows/ear and kernel weight/ear, while at Qlyasan location and
the first generation highly significant mean square due to Genotype were exhibited
for all characters except 300-kernel weight. At Kanipanka location and the second
generation highly significant mean squares due to Genotypes were recorded by
the characters plant height, no. of ears/plant, no. of rows/ear and 300-kernel
weight, while at Qlyasan location and the second generation the mean squares due
to Genotypes were highly significant for all characters. The mean squares due to
parents at Kanipanka location for F1 generation were highly significant for no. of
ears/plant, no. of rows/ear and kernel weight/ear, while at Qlyasan location it was
highly significant for plant height, no. of ears/plant, and kernels weight/row and
it was significant for no. of rows/ear, no. of kernels/row and kernel yield/plant.
At Kanipanka location for second generation the mean squares due to Parents were
highly significant for only 300-kernel weight and significant for plant height, while
at Qlyasan location for the second generation it was highly significant for plant
height and no. of kernels/row, and it was significant for kernels weight/row and
300-kernel weight. The mean squares due to the crosses for the first generation at
Kanipanka location were highly significant for no. of ears/plant, no. of rows/ear and
kernel weight/ear, while it was significant for plant height only. At Qlyasan location
the mean squares due to crosses for the first generation was highly significant for
plant height, kernels weight/row and kernel weight/ear. Regarding to the second
generation at Kanipanka location the mean squares due to the Crosses were highly
significant for 300-kernel weight only. At Qlyasan location it was highly significant
for plant height, no. of kernels/row, kernels weight/row and kernel yield/plant
and significant for no. of rows/ear and 300-kernel weight. The mean squares

8
due to parents versus crosses at Kanipanka location for F1 generation were highly
significant for the characters plant height, no. of ears/plant, no. of rows/ ear, kernel
weight/ear and kernel yield/plant. At Qlyasan location it was highly significant
for all characters except 300-kernels weight. At Kanipanka location for the second
generation the mean squares due to parents versus crosses were highly significant for
all characters except the character kernels weight/row which was significant only. At
Qlyasan location for the second generation it was highly significant for all characters.
The mean squares due to line at Kanipanka location for the first generation were
highly significant for plant height, no. of kernels /row, kernels weight/row, kernel
weight/ear, 300-kernel weight and kernel yield/plant, while at Qlyasan location it
was highly significant for kernel weight/ear, 300-kernel weight and kernel yield/plant
and it was significant for only plant height. At the second generation for Kanipanka
location the mean squares due to line were highly significant for all characters except
300-kernel weight which was significant only. At Qlyasan location it was highly
significant for no. of ears/plant, kernel weight/ear and 300-kernels weight. Regard
to the mean squares due to testers at Kanipanka location for F1 generation it was
highly significant for plant height, no. of ears/plant, no. of kernels/row, kernel
weight/ear, 300-kernel weight and kernel yield/plant, but at Qlyasan location it was
highly significant for no. of kernels/row, kernels weight/row, kernel weight/ear,
300-kernel weight and kernel yield/ plant. For the second generation at Kanipanka
location the mean squares due to testers were highly significant for all characters,
while at Qlyasan location it was highly significant for no. of ears/plant, no. of
rows/ear , no. of kernels/row and kernel weight/ear, but significant only for 300-
kernel weight. The mean squares due to lines ¤ testers at Kanipanka location for
F1 generation were highly significant for plant height, kernel weight/ear, 300-kernel
weight and kernel yield/plant and significant for no. of kernels/row and kernels
weight/row, while at Qlyasan location it was highly significant for all characters
except no. of ears/plant and no. of rows/ear. For F2 generation at Kanipanka
location it was highly significant for all characters except no. of ears/plant which
was significant only, while at Qlyasan location it was highly significant for no. of
ears/plant, no. of kernels/row, kernels weight/row, kernel weight/ear, 300-kernel
weight and kernel yield/plant.
Data in Table 4 represents the average of genotypes at both locations and both
F1 and F2 generations. Parent 1 produced maximum values for the character 300-
kernel weight at both locations and both generations with 85.126, 84.416, 74.312
and 67.524 g. respectively. Parent 2 produced maximum value for the character
no. of ears/plant and kernel weight/row with 1.685 and 8.469 at Kanipanka location
due to F1 generation, while at Qlyasan location produced maximum value for no.
of rows/ear with 18.525. At Qlyasan location for the second generation it gave the
highest value for no. of ears/plant and kernel yield/plant with 0.933 and 87.627g
respectively. Parent 3 gave the highest value for the characters no. of ears/plant
and kernel weight/row with 0.858 and 5.446 at Kanipanka location for the second
generation, while at Qlyasan location for the same generation it produced maximum
value for the characters no. of rows/ear and no. of kernels/row with 13.532 and
26.328. Parent 4 exhibited the highest value for the character kernel weight/ ear
at both locations for the first generation with 94.332 and 132.487 g. respectively,
while at Qlyasan location for the second generation it produced the highest value

9
 for the characters kernel weight/row and kernel weight/ear with 5.974 and 88.052
g. respectively. Regard to parental values, parent 5 gave maximum value for the
characters plant height, no. of rows/ear, no. of kernels/row and kernel yield/plant
with 191.916, 15.669, 32.784 and 124.287 at Kanipanka location for F1 generation,
while at Qlyasan location it gave maximum value for plant height, no. of kernels/row
and kernel yield/plant with 200.454, 40.364 and 181.665 respectively. For the
second generation at Kanipanka location parent 5 also gave maximum value for
the characters plant height, no. of rows/ear, no. of kernels/row, kernel weight/ear,
300-kernel weight and kernel yield/plant with 140.064, 10.939, 19.240, 64.866 and
64.398 g. respectively, while for the second generation at Qlyasan location. Parent
5 produced maximum value for only plant height with 140.022 cm. Regard to
the crosses values, the cross 1¤3 produced maximum value for the character 300-
kernel weight with 95.216 for the first generation at Kanipanka location and for
no. of ears/plant at Qlyasan location for the first generation with 1.739 and the
characters no. of ears/plant, no. of kernels/row, kernel weight/row and 300-kernels
weight with 1.800, 35.579, 10.084 and 92.169 at Qlyasan location due to the second
generation. The cross 1¤5 gave the highest value for the character no. of rows/ear
with 18.536 at Kanipanka location due to the first generation, while at Qlyasan
location for the same generation it exhibited the highest value for plant height, no. of
ears/plant , kernel weight/row and kernel weight/ear with 213.838, 1.739, 12.232 and
192.206 respectively, while for the second generation at Kanipanka location it gave
the highest value for 300-kernel weight with 92.286 g., and at Qlyasan location for
the second generation showed the highest value for no. of rows/ear with 17.648. The
cross 2¤3 produced the highest value for the characters plant height, no. of ears/plant
and no. of kernels/ row with 233.774, 2.641 and 35.276 at Kanipanka location for the
first generation, while at Qlyasan location for second generation it gave the highest
value no. of ears/plant with 6.757. The cross 2¤4 produced maximum values for
the characters 300-kernel weight and kernel yield/plant with 94.497 and 214.341.
At Qlyasan location for first generation, while at Kanipanka location the second
generation the cross 2¤4 gave maximum values for the characters plant height, no.
of ears/plant, no. of kernels/row, kernel weight/row, kernel weight/ear and kernel
yield/plant with 196.536, 1.710, 29.562, 8.289, 113.617 and 205.384 respectively,
while at Qlyasan location for second generation it gave the highest value for the
character plant height with 196.404. The cross 2¤5 showed the highest value for
kernel weight/row, kernel weight/ear and kernel yield/plant with 8.312, 125.032 and
187.149 respectively at Kanipanka location for the first generation, while at Qlyasan
location for the same generation it showed the highest value for no. of rows/ear and
no. of kernels/row with 19.224 and 44.107 respectively. For the second generation
at Kanipanka location it produced maximum value for no. of rows/ ear with 16.380,
while at Qlyasan location for the second generation it gave the superior value for
kernel weight/row and kernel yield/plant with 159.188 and 194.787 g. respectively.
These results are in agreement with previous studies [6; 9; 24].
Table 5 showed the heterosis values for the crosses at both locations for both
generations. For the first generation at Kanipanka location, the cross 2×3 gave
the maximum heterosis values for the characters plant height, no. of ears/plant ,
no. of kernels/row and kernel weight/row with 29.656, 65.245, 41.543 and 4.489
respectively, while the cross 1×3 gave the maximum heterosis values for the char-

10
acters no. of rows/ear and 300-kernel weight with 62.271 and 16.567 and the same
cross obtained maximum negative heterosis value (-13.898) for the character kernel
weight/ row. The cross 2×5 obtained the maximum value with 43.603 for kernel
weight/ear, and the cross 1×4 recorded the maximum value for kernel yield/plant
with 70.963. Concerning Qlyasan location the cross 2×4 gave the heterosis value
for the characters plant height, 300-kernel weight and kernel yield/plant with 10.506,
25.489 and 68.538 respectively, and the cross 1×3 recorded the highest values due
to no. of rows/ear , no. of kernels/row and kernel weight/row with 45.378, 30.600
and 41.608 respectively. Highest value for no. of ears/plant recorded by the cross
1×5 with 63.440, and for kernel weight/ ear the cross 2×5 recorded the highest
value with 70.602. For the second generation at Kanipanka location, The cross
2×4 recorded the highest heterosis values for the characters plant height, no. of
ears/plant , no. of kernels/row, kernel weight/row, kernel weight/ear and kernel
yield/plant with 45.045, 103.168, 188.864, 120.647, 188.856 and 423.853 respec-
tively, while for no. of rows/ear the cross 2×5 recorded the maximum value with
65.070 and for 300-kernel weight the cross 1×5 recorded the maximum value with
34.782. Regarding Qlyasan location the cross 1×3 recorded the highest heterosis
value for the characters no. of ears/plant, no. of kernels/row, kernel weight/row and
300-kernel weight with 104.391, 66.619, 95.370 and 42.752 respectively, and the
cross 1×4 gave the highest value for kernel weight/ear and kernel yield/plant with
95.293 and 154.101respectively. The crosses 2×4, 2×3 and 1×5 gave the maximum
values for plant height, no. of ears/plant and no. of rows/ear with 43.425, 63.486
and 17.648 respectively. The positive heterosis value ratify the over dominance gene
effect for the parent with high value, while the negative heterosis value signify the
partial dominance gene effect of the parent with low value. Positive and negative
heterosis values results agree with [12].
Table 6 represent the estimation of GCA for the parents for both generations
ant both locations. Regarding to the first generation, at Kanipanka location parent 1
exhibited maximum negative GCA value for the characters no. of kernels/row, kernel
weight/ row and kernel yield/plant with -3.672, -0.493 and -11.724 respectively.
Parent 2 showed maximum positive GCA value for the characters no. of kernels/row,
kernel weight/row and kernel yield/plant with 3.672, 0.493 and 11.724 respectively,
while it obtained the maximum negative vale for 300-kernels weight with -3.377.
The highest positive GCA value was obtained by parent 3 for the characters no. of
ears/plant and 300-kernel weight with 0.268 and 4.041, the same parent maximum
negative value -12.171 for kernel weight/ear. Maximum negative GCA value for no.
of rows/plant obtained by parent 4 with -0.410, while parent 5 obtained maximum
negative GCA value for plant height and no. of ears/plant with -7.547 and -0.171
respectively, also recorded positive GCA value for no. of rows/plant and kernel
weight/ear with 0.417 and 13.687 respectively. Concerning to Qlyasan location,
parent 1 gave the maximum negative GCA value with -0.598 and -5.773 for no.
of rows/plant and kernel yield/plant respectively, while the highest positive value
recorded by the same parent for plant height with 1.385. Parent 2 gave the highest
negative value and the highest positive value for plant height and no. of rows/plant
with -1.385 and 0.594 respectively. Parent 3 gave the maximum positive GCA value
for no. of ears/plant and kernel yield/plant with 0.041 and 8.532 respectively, while
it gave the maximum negative value -0.266 for kernel weight/row. The highest

11
 negative GCA value obtained by parent 4 for the characters no. of ears/plant, no. of
kernels/row and kernel weight/ear with -0.046, -1.293 and -24.213 respectively, and
gave the highest positive value for 300-kernel weight with 4.450. Parent 5 exhibited
the highest positive values for no. of kernels/row, kernel weight/ row and kernel
weight/ ear with 2.310, 0.679 and 30.133 respectively, and gave the highest negative
value -4.564 for 300-kernels weight. Regard to F2 generation, at Kanipanka location
parent 4 gave maximum positive GCA effect value for the characters plant height,
no. of rows/plant, no. of kernels/row and kernel yield/plant with 1.648, 0.735,
2.148 and 27.188 respectively, while parent 5 showed the highest GCA effect value
for kernel weight/ear and 300-kernels weight with 13.798 and 1.970 respectively,
and for the character no. of ears/plant produced by parent 3 with 0.073. maximum
negative GCA effect value for the characters plant height, no. of rows/plant, no.
of kernels/row, kernel weight/row and kernel weight/ear produced by parent 3 with
-1.335, -0.901, -2.107, -0.753 and -16.810 respectively, and parent 5 showed the
highest negative values for no. of ears/plant and kernel yield/plant with -0.137 and
-16.989 respectively, and for 300-kernels weight it was -1.925 produced by parent 4.
Respect to Qlyasan location and F2 generation, the highest GCA effect value for the
characters no. of ears/plant and kernel weight/row exhibited by parent 3 with 0.070
and 0.341 respectively, while for the characters plant height and no. of kernels/row
produced by parent for with 2.215 and 2.290 respectively. For the characters no.
of rows/plant and kernel weight/ear produced by parent 5 with 1.803 and 9.887
respectively. Maximum negative GCA effect value for the characters plant height,
no. of rows/plant, kernel weight/ear and kernel yield/plant were -3.223, -1.519,
-9.219 and -6.424 recorded by parent 3, and for the characters no. of ears/plant,
no. of kernels/ row and kernel weight/row were -0.939. -2.027 and -0.412 recorded
by parent 5. The high positive GCA effect values for the parents indicate their
contributions in the inheritance of this character to their hybrids, while the negative
values indicate to the contribution of these parents in reducing the character mean
in some of their hybrids. These results indicated the suitability of hybridization
method to improve these characters. Other studies showed the same results [2; 3;
13; 16].
The estimation of SCA effect values at both locations and both generations repre-
sent in Table 7. Regard to the first generation at Kanipanka location, the cross 1×3
produced maximum positive SCA effect value for 300-kernels weight and maximum
negative SCA effect value for plant height with 6.441 and -7.930 respectively. The
cross 1×4 showed maximum positive SCA effect for the characters no. of ears/plant,
no. of kernels/row and kernel weight/row with 0.085, 0.903 and 0.396 respectively.
The cross 1×5 produced the highest positive value for no. of rows/ear with 0.241
and the highest negative value for kernel weight/ear and kernel yield/plant with -
7.129 and -19.323 respectively. The cross 2×3 showed maximum positive values
for plant height and maximum negative value for 300-kernels weight with 7.913
and -6.441 respectively. The cross 2×4 exhibited maximum negative value due to
SCA effect for the characters no. of ears/plant , no. of kernels/row and kernel
weight/row with -0.085, -0.903 and -0.396 respectively. The cross 2×5 showed the
highest negative value for No. of rows/ear with -0.241 and maximum positive value
for kernel weight/ear and kernel yield/ plant with 7.129 and 19.322 respectively.
At Qlyasan location for F1 generation the cross 1×4 produced maximum nega-

12
tive SCA effect value for most characters including plant height, no. of ears/plant
, kernel weight/row, kernel weight/ear, 300-kernels weight and kernel yield/plant
with -2.879, -0.062, -0.984, -35.869, -7.706 and -22.919 respectively, while max-
imum positive SCA effect values for most characters recorded by the cross 2×4
including plant height, no. of ears/plant , kernel weight/row, kernel weight/ear,
300-kernels weight and kernel yield/plant with 0.879, 0.062, 0.984, 35.869, 7.706
and 22.919 respectively. At Kanipanka location for second generation the cross
1×4 showed maximum negative SCA effect value for the characters plant height,
kernel weight/ear and kernel yield/plant with -3.497, -13.044 and -38.328 respec-
tively, while maximum positive value for the characters no. of kernels/row, kernel
weight/row and 300-kernels weight produced by the cross 1×5 with -3.361, -1.475
and -3.871 respectively. Respect to the second generation at Qlyasan location the
cross 1×3 showed the maximum positive values for no. of ears/plant, no. of
rows/ear, kernel weight/row and 300-kernel weight with 0.178, 0.283, 1.708 and
4.141 respectively. The cross 1×4 showed maximum negative SCA effect value
for plant height with -2.523. The cross 1×5 gave the highest negative SCA for
no. of kernels/row, kernel weight/ear and kernel yield/plant with -5.438, -16.107
and -27.463 respectively. The cross 2×3 showed maximum positive value for no.
of ears/plant with 1.557 and maximum negative value for no. of rows/ear, kernel
weight/ row and 300-kernels weight with -0.283, -1.708 and -4.141 respectively.
Maximum positive value due to SCA effect for no. of kernels/row, kernel weight/
ear and kernel yield/plant recorded by the cross 2×5 with 5.438, 16.107 and 27.463
respectively. The positive SCA effect value caused the increase in the character
value in compare to their parents, while the negative effective of SCA value caused
the decrease in the character value in compare to their parents. Similar results were
recorded by [7; 11; 17; 20; 21].
The proportional contribution of lines, testers and their interactions to total vari-
ance at both locations and both generations represent in Table 8. Respect to the
first generation, at Kanipanka location the testers contributed more to the total sum
square for plant height, no. of ears/plant, no. of rows/ear and kernel weight/row with
75.317, 87.660, 76.625and 76.625% respectively, while lines brought much varia-
tion in the expression of the characters no. of kernels / row, kernel weight/ear and
300-kernels weight with 79.959, 79.959 and 61.536% respectively. Contribution of
line × tester was slightly greater than that of testers and lines for kernel yield/plant
with 50.212%. At Qlyasan location, line × tester interactions contributed more to
the total sum square for all the studied characters except no. of rows/ear and kernel
weight/row each with 70.319%. Concerning the second generation, at Kanipanka
location testers contributed more to the total sum square for no. of ears/plant, no.
of rows/ear and kernel weight/row with 59.374, 40.546 and 40.546% respectively,
while the interaction of line × tester brought much variation in the expression of
plant height (73.808%), no. of kernels/row (46.721%), kernel weight/ear (46.721%),
300-kernels weight (67.298%) and kernel yield/plant (54.128%). At Qlyasan loca-
tion the testers brought much variation in the expression of the characters plant
height, no. of rows/ear and kernel weight/row with 60.057, 91.228 and 91.228%
respectively. Line × tester interactions contributed more to the total sum square
for no. of ears/plant, no. of kernels/row, kernel weight/ear, 300-kernels weight and
kernel yield/plant with 75.895, 85.634, 85.634, 93.352 and 90.502% respectively.

13
 The results of the present study revealed large variation between parents and hybrids
for the characters under study. Compared to the parents, hybrids were earlier, taller,
bearing more effective tillers, and had higher kernel yielding.

Conclusion

The crosses 2×4 and 2×5 possess the best SCA and gave the desirable positive
heterosis values for kernel yield and some of its components. These crosses are
great promising to be utilizing in breeding programs of maize in Iraqi Kurdistan
Region, which possess genetic factors for yield potentiality.

14
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20
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23
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References

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Tamil Nadu Agricultural University, Undergraduate Program. COMBATORE
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Maize (Zea mays L.)’, in, College of Agri. , Univ. of Baghdad, (In Arabic).
[3 ] Al- Falahy, A. A. [2002], ’Genetic Parameters in Diallel cross in Maize’, in,
College of Agri, Univ. of Baghdad. (In Arabic).
[4 ] Al-Barodi, M. M. M. [1999], ’Partial diallel for inbred line of maize’, in,
College of Agriculture, Baghdad University. (In Arabic).
[5 ] Ali, G., Ishfaq, A., Rather, A., Wani, S. and Zargar, M. [2006], ’Combining
ability analysis for grain yield and maturity traits in high altitude maize inbreds
(Zea mays L.)’. National Journal of Plant Improvement, 8(2), 128-131.
[6 ] Barakat, A. A. and Ibrahim, M. H. A. [2006], ’Heterosis and combining
ability in yellow Maize’. Mansoura University, J. Agri. Sci, 31(8), 4849-4860.
[7 ] Desai, S. A. and Singh, M. [2001], ’Combining ability studies for some mor-
phological, Physiological and biochemical traits related to drought tolerance in
maize’. Indian J. Genetic. and plant Breeding, 61, 34-36.
[8 ] Dillen, B. S. [1975], ’Application of partiall diallel crosses in plant breeding’.
A review, crop improvement, 2, 1-7.
[9 ] Flint-Garcia, S. A., Buckler, E. S., Tiffin, P., Ersoz, E. and Springer, N. M.
[2009], ’Heterosis is prevalent for multiple traits in diverse maize germplasm’.
PloS one, 4(10), e7433.
[10 ] Gautam, A. S. [2003], ’Combining ability studies for grain yield and other
traits in inbred lines of maize’. Crop Res., 26(3), 482-485.
[11 ] Hallauer, A. [1999], ’Heterosis What Have We Learned? What Have We
Done? Where Are We Headed?’. The Genetics and Exploitation of Heterosis
in Crops(thegeneticsande), 483-492.
[12 ] Iqbal, M., Khan, K., Rahman, H., Khalil, I., Sher, H. and Bakht, J. [2010],
’Heterosis for morphological traits in subtropical maize (Zea mays L)’. May-
dica, 55(1), 41.
[13 ] Kara, S. M. [2001], ’Evaluation of Yield and Yield Components in Inbred
Maize Lines I. Heterosis and Line x Tester Analysis of Combining Ability’.
Turkish Journal of Agriculture and Forestry, 25(6), 383-391.
[14 ] Kempthorne, O. [1957], ’An introduction to genetic statistics’.
[15 ] Keskin, B., Yilmaz, I. H. and Arvas, O. [2005], ’Determination of some
yield characters of grain corn in eastern Anatolia region of Turkey’. Journal
of Agronomy.

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 [16 ] Malik, S. I., Malik, H., Minhas, N. M. and Munir, M. [2004], ’General and
specific combining ability studies in maize diallel crosses’. Int. J. Agri. Biol,
6(5), 856-859.
[17 ] Mohammad, S. M. [2005], ’Reciprocal Diallel Crosses and estimation of
Genetic Component of Maize (Zea mays L.) in Kurdistan’, in, College of Agri.
Univ. of Salahaddin. (In Arabic).
[18 ] Mohammed, I. M. and Al- Jubouri, J. M. A. [2007], ’Estimation some genetic
parameter by (line ¤ tester) analysis of corn (Zea mays L.)’. Kerkuk Journal
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[19 ] Rezaei, A., Yazdisamadi, B., Zali, A., Tallei, A., Zeinali, H., Rezaei, A., Voll-
mann, J., Grausgruber, H. and Ruckenbauer, P. [2004], ’Estimate of heterosis
and combining ability in maize (Zea mays L.) using diallel crossing method’,
in, Genetic variation for plant breeding. Proceedings of the 17th EUCARPIA
General Congress, Tulln, Austria, 8-11 September 2004., BOUK-University of
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[20 ] Sadalla, H. A., Mustafa N. R. and A., K. S. [2012], ’Pollen Source Effect on
Maize kernel Oil Concentration’. University of Salahaddin, Hawler, Zanko J.
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[21 ] Sharief, A., El-Kalla, S., Gado, H. and Abo-Yousef, H. [2009], ’Heterosis in
yellow maize’. Australian Journal of Crop Science, 3(3), 146.
[22 ] Sharma S., Narwal, M. S. and Kumar, R., Dass S. [2004], ’Line ¤ tester
analysis in maize (Zea mays L.)’. Forage Res, 30, 28-30.
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tative genetic analysis’. Biometrical methods in quantitative genetic analysis.
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through Line¤ Tester analysis in maize (Zea mays L.)’. Plant Archives, 11(1),
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[25 ] Vasic, N., Ivanovic, M., Peternelli, L., Jockovic, D., Stojakovic, M. and
Bocanski, J. [2001], ’Genetic relationships between grain yield and yield com-
ponents in a synthetic maize population and their implications in selection’.
Acta agronomica hungarica, 49(4), 337-342.

26
 Advances in Detecting Breast Mass by 2D Ultrasound
Heamn Noori Abduljabbar and Sardar Yaba Perxdr
Department of Physics, College of Education Shaqlawa, Salahaddin University,
Kurdistan-Iraq
[email protected]
Abstract
Breast mass can be classified to cystic and solid masses and in both there are
many sub types. The aim of this study is how to improve the early detection of
breast masses (cystic/solid) by 2D ultrasound scanning, by running the ultrasound
images of the abnormal cases in image processing software by applying some options
such as image enhancement, edge detection, and filtration, furthermore comparing
between the resulted 2D images by 3D images, and Elastography of breast mass
images. 783 patients had been scanned in Harer hospital by 2D ultrasound B-
Mode Linear probe with 7.2 MHz frequency, ( 47) of them diagnosed as breast
mass, the images saved then applied to image processing program. The result was
the obtained 2D images are much more clear for diagnosis after running in to image
processing, the mass can easily be detected and classified in to the correct type
without using 3D high cost ultrasound. .
Keywords: Breast mass, Cyst, Fibroid, Image enhancement, Edge detection

Introduction

Complex cystic masses can defined as lesions composed of echogenic (solid) compo-
nents and anechoic (cystic), while in complicated cysts, the echogenic fluid content
of which imitates a solid lesion. Complex breast masses are requiring histological
verification by surgical ablation and/or percutaneous biopsy, then it may be benign
or high risk (hematoma, an abscess, papilloma, fibrocystic mastopathy, fat necro-
sis, a phyllodes tumor) as much as malignant (necrotic cancer, papillary cancer, a
ductal carcinoma in situ, metastases). In order not to fail to recognize high risk or
malignant lesions it is requiring appropriate management Histopathological corre-
lation is essential to ensure that the samples are representative and concur with the
ultrasound appearance. The biopsy technique must be adapted to each case and it
is often necessary to insert a coil during the procedure [6,9].

Cysts are by far the most frequently encountered breast condition and they are
usually asymptomatic. We can describe the cystic lesions according to the BIRADS
lexicon [1,2] typically to benign, simple cysts which is imperceptible wall, strictly
anechoic content, posterior acoustic enhancement; probably benign, while compli-
cated cysts are imperceptible wall, with or without posterior enhancement, homo-
geneous hypoechoic content or with more or less sloping fluid/fluid levels; complex
cysts of indeterminate nature association of cystic and solid components, thick wall
or internal septum/septa. The challenge for the radiologist is differentiating between
a complicated cyst and a complex cystic mass [5].

27
 Anatomy

The breast generally refers to the front of the chest and medically specifically to
the mammary gland. (The word ”mammary” comes from ”mamma,” the Greek and
Latin word for the breast, The mammary gland is a milk-producing structure that is
composed largely of fat cells (cells capable of storing fat). The fat deposits are laid
down in the breast under the influence of the female hormone estrogen. It is consist
of three layers mainly which are fatty layer, lactic (nodules), and muscular layer
which makes the base wall of the breast [7].as shown in figure 1. Breast Imaging

Figure 1: Breast anatomy

modalities

1. Magnetic Resonance Imaging MRI

2. Computed Tomography CT scan
3. Ultrasound US
4. Nuclear Medicine
5. Mammography
6. Thermography

Ultrasound

Ultrasound is used extensively to investigate the nonlactating breast to distinguish

between abnormalities (e.g. cysts, fibroadenomas and malignant changes) and nor-
mal structures (e.g. milk ducts, Cooper’s ligaments, lymph nodes, glandular and
adipose tissue) [4]. In addition, recent technological advances have increased the
ability of ultrasound to resolve very fine structures such as ducts as small as 0.5
mm [4].

28
Examination techniques

Before performing and interpreting a breast ultrasound examination, as for any

imaging examination, it is essential to have available all the clinical and radiological
data. (B mode) should meet the following criteria, linear probes of frequencies
between 7.2 and 20 MHz. The frequency must be adapted to the breast volume
and the location of the lesion in the breast (lower frequency for deep lesions, high
frequency for superficial lesions); total gain and TGC should be used to obtain an
adequate and homogeneous signal from the skin to the depth of the pectoral muscle;
the lesion must be analyzed in at least two planes: transverse and longitudinal. The
use of additional modes helps characterize the lesion; in case of using advanced
or developed machines, while by using 2D low cost machine there is only one
mode. Images to be saved then processed by image processing software. Image
enhancement with fuzzy method of edge detection are used in this study.

Ultrasound signs of complicated cysts

Complicated cysts show all the aspects of simple cysts except for the content, which
is finely echogenic. They may have a fluid level or internal echoes that correspond
to debris and which are displaced slowly with changes in the patient’s position[5].

Figure 2: Show types of breast mass

The research presents an approach for improving range image segmentation,

based on fuzzy regularization of the detected edge. In this paper, we present a
way to define the concept of fuzzy image segmentation, which has not been clearly
defined in the literature. In this work the term Fuzzy image segmentation is char-
acterized by means of a fuzzy set over the set of edges, which can be understood
as the fuzzy boundary of the image. Also we discuss a visualization of an image
segmentation in terms of edge detection. But first, we define two concepts of crisp
image segmentation on an image network, one based on nodes and the other one
focusing on edges.

29
 This study focuses on fuzzy logic based edge detection in smooth and noisy
clinical images. The proposed method (in noisy images) employs a 3 × 3 mask
guided by fuzzy rule set. Moreover, in case of smooth clinical images, an extra
mask of contrast adjustment is integrated with edge detection mask to intensify the
smooth images.

Fuzzy approach for edge detection Image processing possesses vagueness and
ambiguity and fuzzy deals with the data uncertainty [8]. Fuzzy logic provides math-
ematical framework. General procedure of fuzzy logic is shown in Fig. 3. Fuzzy
methodology is generated to deal with brightness, edges and geometric features.
Fuzzy logic also deals with subjective concepts. There are different possibilities for
edge detection based on fuzzy logic.In case of edge detection uncertainty is high
in neighbourhood pixel. One method is to define a membership function and using
IF-THEN rules general edge detections can be performed. Membership function
is determined heuristically. It is fast but has some limitations. Fuzzy clustering,
fuzzy rule based system, fuzzy set theory are the methods of edge detection based
on fuzzy logic.

Theory

Fuzzy Logic edge detection

As known, the fuzzy has three steps. Two-dimension Gradient Gx and Gy com-
puted and fuzzified by two Gaussian membership functions. The rules that modify
membership values is to be white color for both two gradient with zero value else
black color if one of them don’t zero. Lastly, the resulted defuzzified by triangular
membership function. Table 1 shows fuzzy sets for input and output variables.

Figure 3: Basic block diagram

Table 1: Fuzzy sets for input and output variables.

Fuzzy sets Range MF type
P 1input [0.10] Gaussian
P 2input [0.10] Gaussian
White [0.1 1 1]
Poutput Triangular
Black [0 0 0.7]

30
Results and Discussion

Simple Breast cyst

Simple breast cyst defined in the 2D ultrasound in the original image as shown in
figure 4 having thin regular wall with clear fluid filled while figure 5 shows the same
image after using image enhancement teqnique in image processing which improves
the image quality as or close to 3D ultrasound image. The same image after

Figure 4: Simple breast cyst (original 2D image)

Figure 5: Simple breast cyst after making image enhancement

using edge detection as shown in figure 6, by using Fuzz edge detection way, the
cyst is more clear and can be define even better than infra-red IR, or other imaging
modalities, also the image quality is improved and looks much more clear such as
3D, or 4D ultrasound. As shown in (figure 6)

Complex breast cyst

Complex breast cyst is another type of breast cysts it can be recognized by ultrasound
scanning as it has thick irregular outline and acoustic shadows as shown in figure 7
also the fluid inside it is thick fluid not clear echogenicity. While figure 8 shows
the same cyst after running the image into image processing software by making

31
 Figure 6: Explains the using of Fuzzy Logic Edge detection way which shows the cyst much more
clear

image enhancement which give us more information about the cyst from figure 7,
moreover we can see the intensity of the wall thickness clearly after using fuzzy
edge detection for the same image it helps also to detect the severity of the cyst by
defining the outline and the acoustic shadow as shown in figure 9 below. It is clear
that the image quality is higher and better after using Fuzzy edge detection to detect
complex breast cyst than Elastography as shown in figure 10.

Figure 7: Shows the cyst with thick irregular outline and has multiple acoustic shadows (original
2D image)

Figure 8: Complex breast cyst shows wall thickness the red arrow and thick fluid the blue arrow
better

32
Figure 9: The wall thickness clearly defines the severity of the cyst and the intensity of the cyst
outline clearly shows the rate of the outline’s solidity and compatibility.

Figure 10: Complex breast cyst by using Elastography technique.

References

[1 ] Neil J. Diffusion imaging concepts for clinicians. J Magn Reson. Imaging

2008;27:1Ű8.
[2 ] American College of Radiology. ACR BI-RADS ultrasound. In: ACRbreast
imaging reporting and data system, breast imaging atlas.Reston, Va: American
College of Radiology; 2003.
[3 ] Venta LA, Kim JP, Pelloski CE, Morrow M. Management of com-plex breast
cysts. AJR Am J Roentgenol 1999;173(5):1331Ů6.
[4 ] Doshi DJ, March DE, Crisi GM, Coughlin BF. Complex cysticbreast masses:
diagnostic approach and imaging-pathologiccorrelation. Radiographics 2007;
27 (Suppl. 1):S53Ů64.
[5 ] Cha JH, Moon WK, Cho N, Chung SY, Park SH, Park JM, et al.Differentiation
of benign from malignant solid breast masses:conventional US versus spatial
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[6 ] Candelaria RP, Hwang L, Bouchard RR, Whitman GJ. Breastultrasound:
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33
 BI-RADS classification with subcate-gories of category 4. Eur Radiol 2013;
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[8 ] Zhi H, Ou B, Luo BM, Feng X, Wen YL, Yang HY. Comparison of Ultrasound
Elastography, Mammography, and Sonography in the Diagnosis of Solid Breast
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[9 ] H.W. Kim, J.M. Chang, W.K. Moon, N. Cho, A. Yi, H.R. Koo, et al. Intra-
ductal mass on breast ultrasound: final outcomes and predictors of malignancy
AJR Am J Roentgenol, 200 (4) (2013), pp. 932Ű937

34
 Antimicrobial Activity of Silver Nanoparticles Synthesized From Banana’s
Peel Extract
Huda Z. Majeed and Rana H. Al-Shammari
College of Science, Al-Mustansiriyah University, Baghdad, Iraq
[email protected]
Abstract
The peels of banana represent less than half of the total weight of fresh banana
and the peels are usually discarded as waste after the inner part is eaten. Processing
industries are suffering from high amounts of peels, in order to benefit from after
transforming into more valuable product. Peels of banana can be used for allergies
and skin irritations at home, for biofuel making, energy related activities, paper,
cosmetics, environmental cleaning, organic fertilizer, processes related to biotech-
nology and in nanotechnology. Multi-drug resistant bacteria is a important problem
in the infectious diseases treatment and the random use of wide-spectrum antibiotics
has produced increased antibiotic resistance for many different bacterial pathogens.
Nanotechnology have opened wide door in the synthesis of nanoparticles. Since
old days, silver was known for its antibacterial effects and it has been used for
prevention and control of many different infections. The present study was aimed
to synthesis of silver nanoparticles from banana peels, which were characterized by
UV-VIS Spectroscopic Analysis and Scanning electron microscope (SEM), which
were done at Nanotechnology Center in UOT in Iraq. The nanoparticles exhibited
highest absorbance peak at 440 nm in UV-Vis spectroscopy. The AgNPs surface
morphology from SEM images reveals formation of well-dispersed Ag-NPs with
diameter between (1- 239) nm, the AgNPs diameter average was (80 nm), and the
presence of silver was confirmed. The AgNPs prepared from banana’s peel extract
had antimicrobial activity with the most effective concentration (30 µg/µl) against
E.coli. and (40 µg/µl) against Salmonella sp., Pseudomonas aeruginosa and Can-
dida albicans.
Keywords: Silver nanoparticles, Banana peels, Multidrug resistant bacteria and An-
timicrobial activity.

Introduction

Nanobiotechnology is a field that combines both biological sciences and nanotech-

nology. It provides a ground for the green synthesize of ecofriendly nanoparticles
with the aid of biological sources like plants and microorganisms [1]. The plants
or plants extract, which act as reducing agents for nanoparticles synthesis, are more
valuable than other biological materials [2], because they exclude the process of
culturing and preserving of the cell, and can scaled up for large-scale nanoparticle
synthesis [3]. Moreover, plant-mediated nanoparticles synthesis is benefit because
it is low-cost, ecofriendly, a single-step process for biosynthesis method and safe
for human [4]. Different part materials of plants could be used for synthesis such
as extracts [5], fruit [6], bark [7], fruit peels [8]. The study aimed to prepare silver
nanoparticles (AgNPs) from peels of banana, characterizing these particles by UV-
VIS Spectroscopic and Scanning electron microscope (SEM) . After that, Detecting

35
 of Antimicrobial activity against multi-drug resistant microbes was done.

Materials and Methods

Banana peel Extract preparation

Twenty gram of fresh banana peel was added into 100 ml of distilled water and
boiled for 10 minutes to get the banana peel extract, then cooled and filtered using
Millipore filter (0.22) and filtrate was used for the synthesis of AgNPs [9].

Silver Nanoparticle synthesis

10 ml of freshly prepared filtrate was added to 90 ml of 1mM silver nitrate solution to

reduce Ag+ ions and incubated at room temperature .The color change was observed
, then centrifuged at 10,000 rpm for 5 min and pellets were collected, dried and
stored [9].

Silver nanoparticles characterization

UV- visible spectroscopy

A UV visible absorption spectrophotometer with a resolution of 1.0 abs between

300 to 550 nm was used in order to get UV visible spectroscopy analysis in aqueous
solution [10,11].

Scanning Electron Microscope (SEM) analysis

Scanning Electron microscopy used for morphological characterization .Sample was

made by putting the drop of sample suspension over carbon coated grid then it was
dried, examined and photographed in SEM [12].

Antimicrbial activity of silvernanoparticles

Wells were made on the agar plates using a gel puncture to about 10 mm diameter
in Nutrient agar medium. Strains were swabbed uniformly onto the individual plates
using sterile cotton swabs. 100 /µ g of dried AgNPs were dispersed in 100 /µ l of
distilled water. The dispersed solution was impregnated in the well at (25, 30, 35
and 40) /µ l, to get the concentration of (25, 30, 35 and 40) /µ g/µ l respectively, to
perform the well diffusion method against bacterial isolates which include (E.coli ,
Ps. aeruginosa , Salmonella , Staphylococcus aureus and Bacillus subtilis). Sterile
D.W. is used as control. After 24 h, the diameters of inhibition zones around
the wells were measured in millimeter. The size of the circular inhibition zone is
directly proportional to the antimicrobial effect of the biosynthesized AgNPs against
microbial pathogens [13].

36
Results and Discussion

Synthesis of Silver Nanoparticle from banana peel

When extract of banana peels incubated with aqueous solution of silver nitrate, the
color changed from pale brown to dark brown color due to the reduction of Ag+
ions, suggesting the formation of AgNPs as shown in Figure (1).

Figure 1: Synthesis of Silver Nanoparticle from banana peel extract (A)1 mM of AgNO3 control
(B) 1 mM of AgNO3 after 24hrs (B) mixing with Banana peel extract

UV-Vis spectroscopy

UV-Vis Spectroscopy used to examine the mean size of AgNPs in aqueous suspen-
sion. The assurance of formation and stability of AgNPs was detected by using
UV-Vis analysis. The extract with silver nitrate showed the sharp peak around 440
nm with high absorbance which is very specific of AgNPs as shown in Figure (2).

UV-Vis spectroscopy

Scanning Electron Microscope (SEM) analysis SEM was used to examine the mor-
phology of AgNPs. SEM image showed the high density nanoparticle synthesized
by extract of banana peels. It was done at the magnification of 9.7mm were Mono
dispersed, round shaped AgNPs were found. this shape of the particle confirms the
presence of AgNPs.
Banana peels are mainly composed of pectin, cellulose and hemicelluloses and
the functional groups associated with these polymers ,also the proteinaceous matter
may play a role in reducing the silver salt to Ag+ . Biological components are known
to interact with metal salts via these functional groups and mediate their reduction
to AgNPs [14].

37
 Figure 2: UV-Vis spectroscopy of Silver Nanoparticle from banana peel extract

Figure 3: (SEM) analysis of Silver Nanoparticle from banana peel extract

38
Antimicrobial activity of silver nanoparticles

The biologically synthesized AgNPs showed excellent antimicrobial activity against

clinically isolated Multidrug-resistant human bacteria such as Gram-positive bac-
teria Salmonella, Staphylococcus aureus, B. subtilis and Gram-negative bacteria ,
Ps. aeruginosa and E.coli as shown in Figure [5]. The Gram-negative bacteria Ps.
aeroginosa and E. coli showed maximum inhibition which may be due to the cell
wall of Gram + bacteria composed of a thick peptidoglycan layer, which composed
of linear polysaccharide chains cross linked by short peptides, forming more rigid
structure leading to difficult entrance of the AgNPs compared to the Gram-bacteria
where the cell wall had thinner peptidoglycan layer [15]. The high bactericidal ac-
tivity is certainly due to the silver cations released from AgNPs that act as reservoirs
for the Ag bactericidal agent. Big changes in the membrane structure of bacteria
as a result of the interaction with silver cations lead to the increased membrane
permeability of the bacteria [16].

Figure 4: Antibacterial activities of AgNPs prepared from banana peel extract (A) Ps. aeroginosa
(B) Salmonella (C) E. coli (D) Staphylococcus aureus (E) B. subtilis C: control (1 mM
of AgNO3) 1: 25 µ g/µ l 2: 30 µ g/µ l 3: 35 µ g/µ l 4: 40 µ g/µ l

The biocidal properties of AgNPs against microorganisms could be explained by

different effects. First, AgNPs attach to the cell surface which is negatively charged,
change the physical and chemical properties of the cell membranes and the cell
wall and defect important functions such as permeability, osmoregulation, electron
transport and respiration of the bacterial cell [17]. Second, AgNPs can cause further
damage to bacterial cells by penetrating the cell, where they interact directly with
DNA, proteins and other phosphorus- and sulfur-containing cell components [18].
Third, AgNps release silver ions, leading to an amplified biocidal effect, which is
size- and dose-dependent [19].

Conclusions

Banana peel is a waste product in the processing industry could be used for syn-
thesizing AgNPs, which in turn, were produced by the direct exposure of silver
nitrate with peel extract in aqueous media without the intervention of any external
chemicals. Therefore, this reaction pathway satisfies all the conditions of a 100 %

39
 green chemical process. The methodology employed here is very simple, easy to
perform, inexpensive, and eco-friendly.

References

[1 ] Guangquan L, Dan H, Yongqing Q, Buyuan G, Song G and Yan C (2012)

Fungus mediated green synthesis of silver nanoparticles using Aspergillus ter-
reus. Int J Mol Sci 13:466-476
[2 ] Valli, J S and Vaseeharan, B (2012) Biosynthesis of silver nanoparticles by
Cissus quadrangularis extracts. Materials Letters, 82, 171e173.
[3 ] Saxena, A , Tripathi, R M , Zafar, F , and Singh, P (2012) Green synthesis of
silver nanoparticles using aqueous solution of Ficus benghalensis leaf extract
and characterization of their antibacterial activity. Materials Letters, 67, 91e94.
[4 ] Kumar, D A , Palanichamy, V and Roopan, S M (2014) Green synthesis of
silver nanoparticles using Alternanthera dentate leaf extract at room temper-
ature and their antimicrobial activity. Spectrochimica Acta Part A-Molecular
and Biomolecular Spectroscopy, 127, 168e171.
[5 ] MubarakAli, D , Thajuddin, N , Jeganathan, K and Gunasekaran, M (2011)
Plant extract mediated synthesis of silver and gold nanoparticles and its an-
tibacterial activity against clinically isolated pathogens. Colloids and Surfaces,
Biointerfaces, 85: 360-365.
[6 ] Prathna, T C , Chandrasekaran, N , Raichur, A M and Mukherjee, A (2011)
Biomimetic synthesis of silver nanoparticles by Citrus limon (lemon) aque-
ous extract and theoretical prediction of particle size. Colloids and Surfaces,
Biointerfaces, 82: 152-159.
[7 ] Satishkumar, M , Sneha, K , Won, S W, Cho, C W , Kim, S and Yun, Y S
(2009) Cinnamon zeylanicum bark extract and powder mediated green synthesis
of nano-crystalline silver particles and its antibacterial activity. Colloids and
Surfaces, Biointerfaces, 73: 332-338.
[8 ] Bankar, A , Joshi, B , Kumar, A R and Zinjarde, S (2010) Banana peel
extract mediated novel route for the synthesis of silver nanoparticles. Colloids
and Surfaces A-Physicochemical and Engineering Aspects, 368: 58-63.
[9 ] Nirmala, J. and Keerubavathi (2016) Synthesis of silver nanoparticles us-
ing fresh banana peel musa extract. Science & Technology Research. The
International Quarterly journal, 2[5]: 73-82
[10 ] Narasimha G , Praveen B , Mallikarjuna K , and Deva Prasad Raju, B (2011)
Mushrooms (Agaricusbisporus) mediated biosynthesis of silver nanoparticles,
characterization and their antimicrobial activity, Int. J. Nano Dim, 2[1]: 29-36.
[11 ] Haytham ,M M (2015) Green synthesis and characterization of silver nanopar-
ticles using banana peel extract and their antimicrobial activity against repre-
sentative microorganisms. J Radiation Res Appl Sci, 8: 265-275.

40
[12 ] Schaffer ,B, Hohenester U, Trugler A and Hofer F (2009) High-resolution
surface Plasmon imaging of gold nanoparticles by energy-filtered transmission
electron microscopy, Phys Rev B , 79.
[13 ] Prabhu N, Divya Raj, Yamuna gowri K, Ayisha siddiqua S, Joseph puspha
innocent D (2010) Synthesis of silver phyto nanoparticles and their antibacterial
efficacy. Digest J Nanomater Biostruct, 5:185Ű189.
[14 ] Bar H, Bhui DK, Sahoo GP, Sarkar P, Pyne S, Misra A (2009) Green syntheis
of silver nanoparticles using seed extract of Jatropha curcas. Colloids Surf A
Physic- ochem Eng Asp ,348 :212Ű216
[15 ] Shrivastava S, Bera T, Roy A, Singh G, Ramachandrarao P, Dash D (2007)
Characterization of enhanced antibacterial effects of novel silver nanoparticles.
Nanotechnology ,18: 225103Ű225111
[16 ] Sondi I and Salopek-Sondi B (2004) Silver nanoparticles as antimicrobial
agent: a case study on E. coli as a model for Gram-negative bacteria. J
Colloid Interface Sci ,275: 177-182
[17 ] Marambio-Jones, C and Hoek, E M (2010) A review of the antibacterial
effects of silver nanomaterials and potential implications for human health and
the environment. J Nanoparticle Res , 12: 1531-551.
[18 ] Nel, A E , MĂadler, L , Velegol, D , Xia, T , Hoek, E M, Somasundaran,
P (2009) Understanding biophysicochemical interactions at the nano-bio inter-
face. Nature Materials, 8: 543-557.
[19 ] Liu, J Y , Sonshine, D A , Shervani, S and Hurt, R H (2010) Controlled re-
lease of biologically active silver from nanosilver surfaces. ACS Nano, 4:6903-
6913.

41
42
 The Effect of Cutting on Yield and Its Components of Six Genotypes of
Barley Under Slaimani Condition
Shara Jalal Hama, Shang Haseeb Abdulqader, Bestoon Omare and Kazal Rashed
Ahmad
Faculty of Agricultural Sciences, University of Sulaimani, 46001 Sulaimani,
Kurdistan-Iraq
[email protected]
Abstract
A research study was conducted at winter season of 2012-2013 , at Qlyasan Agri-
culture research station to evaluate the effect of cutting treatments on the yield and
its, related components of six-rowed barley (Hordeum vulgare L.) genotypes. Yield
and its components were analyzed according to split plot layout within Completely
Randomized Block Design with 3 replications was used. The 6 genotypes were
implemented in the main plots and the cutting treatments (cutting and non-cutting)
were implemented in the sub plot regards to. Dry forage yield were analyzed ac-
cording to Completely Randomized Block Design with 3 replications. Comparisons
between means were carried out by the lease significant difference (L.S.D) at 5%
level of significance.
Keywords: Barley, Cutting forage, grain yield and yield components.

Introduction

Barley (Hordeum vulgare L) is the forth most- produced cereal after maize, rice
and wheat [1] in the world. It has remained as a main food source in western and
eastern Asia as well as in Himalayan nations and in northern and eastern Africa[2].
Barley serves as a major animal fodder, base malt for beer and certain other distilled
beverages. is one of the main crops of old world agriculture. Remains of barley seeds
found at archeological sites in the Fertile Crescent indicate that about 10000 years
ago the crop was domesticated there from its wild prorated Hordeum spontaneous
[3]. Barley is most tolerance cereal grain species to drought, cold, and salt [1]. It
is worth noting that the toleration of some forage crops to the cutting while keeping
a high productivity of grains is related to the genotype [4, 5 and 6]. There are
different Studies about the effect of varieties and cutting treatment on the grains
and forage production according to the varied genotypes and environment [7, 8 and
9]. Some researchers noted that both genotypes and cutting have negative effect
on the forage and grain yield production. [10, 11, 9 and 12]. The main aim of
the present study is to recognize the effect of genotypes, cutting and non-cutting
treatment on the productivity of forage yield, grain yield and its related components
under sulaimani conditions.

Material and Methods

The present study was carried out at Qlyasan research station, Faculty of Agriculture,
University of Sulaimani (N 35◦ 34’ 307’, E 45◦ 21’ 992’,at altitude of approximately
765 m)during the rabi season from December 2012 to April 2013 to study the effect

43
 of variety and cutting treatments on the yield and its components of barley. The
experimental area plots were ploughed, harrowed and well leveled. Table (1 and 2)
shows the metrological data, physical & chemical properties of the soil study.The
experiment was arranged as split plot layout in RCBD with three replicates. The
six genotypes of barley were implemented in the main plots and the cutting treat-
ments (cutting and non-cutting) were implemented in the sub plots. Each main plot
consisted of two subplots with 4 m width by 1 m length and each sub plot consists
of four rows of 4 m long and 0.25 m between rows. .Seeds cultivated at rate 100
Kg h−1 on Dismember 29, 2012. All plots were fertilized with 60Kg P2 O5/ha as
triple super phosphate, which was broad casted before sowing. Mechanical weed
control was carried out three times to keep the experimental plots free from weeds.
The cutting treatment was done, when plant height was 25-30 cm, and the Cutting
height is usually conducted at 6-8 cm above ground level. The six genotypes of
six-rowed barley (Hordeum vulgare L.) names were:
1. ESP/3/RHODES/CI14100/LINEE527/4/PETUNIA1
CBSS99Y00398T-0TOPM-5Y-1M-1Y-2M0Y
2. ATACO/BERMEJO/HIGO/3/CLN-B/80.5138//Ě
CBSS99M00320T-L-2M-1Y-1M-0Y
3. ALPHA-BAR/DURRA//CORACLE/3/ALELI/4/Ě
CBSS99Y00032S-11Y-2M-2Y-1M-0Y
4. EBC(A)/PALTON//CABUYA
CBSW99WM00073T-AA-1M-1Y-1M-0Y
5. QUINN/ALOE//CARDO/3/CIRU
CBSS99M00038S-11M-1Y-1M-0Y
6. Check variety (Arivate)
Dry forage yield were analyzed according to simple CRBD, Statistical analyses
were conducted using JMP, version 7 [13]. The LSD test was done to find the
significant differences between treatments at 5% probability level .The mature plants
were harvested on 5 of June 2013 for estimating biological yield, seed yield and
yield components.

Result and discussion

Data in Table (3) explain the effect of genotypes on yield and its components char-
acters. It was noticed that this effect was significant for yield and all its components
except straw yield and harvest index. Genotype (2) oultyielded the rest significantly
in most characters. The variation was attributed to the differences in genetic poten-
tial among genotype. This reselt is similar to what obtained by [14 and 15]. Barley
genotypes showed significant differences among them in both biological yield and
1000 Űgrain weight characters similar results recorded by [16]. But [17] showed
that the variation among barley genotypes for the number of spikes m-2 and grain
yield characteristic was significant .Data presented by [17] showed remarkable differ-
ences among barley genotypes in their spick weight ,either in single or in combined

44
 Table 1: Some physical and chemical properties of soil at Qlyassan
Soil properties Values
Soil texture (P.S.D) Clay
Sand ( g.kg−1 ) 41.0
Silt ( g.kg−1 ) 430.50
Clay ( g.kg−1 ) 528.50
E.C. ( dS.m−1 ) 0.61
pH 7.32
O.M. ( g.kg−1 ) 21.60
CaCO3 ( g.kg−1 ) 107.0
Total N (m g.kg−1 ) 1.07
K+ 0.32
+
Na 0.41
Soluble cat ions and anions (mmolesL−1 )
Ca+ + 1.49
Chloride (CL− ) 0.57

Table 2: Average air temperature and rainfall during the season of 2013-2014 at Qlyasan location
Qlyasan
Average o
Month temperature C Rain fall mm
MAX MIN 2012-2013
October 35.6 10.1 27.5
November 28.0 5.8 137.9
December 18.3 1.4 104.9
January 16.1 -1.8 171.5
February 19.3 1.2 57.2
March 27.4 2.8 16.7
April 31.7 7.5 14.0
May 37.5 12.0 43.0
June 40.4 23.4 0.0
Total 572.7

analysis of variance. This agrees with what has been recorded [7, 5, and 9].
Data represented in table (4) explain the effect of cutting treatments in yields and
yield components characters .The non cutting treatment predominated the cutting
treatment significantly for all characters with the exception of average spick weight
only with 14918.47, 7255.18, 473.61, 5747.24, 9171.81, 0.39 and 38.80 for the
characters Biological yield, Spike yield, No. of Spike m2 , grain yield, straw yield,
harvest index and 1000- grain weight respectively. This agrees with what has been
recorded by [10 and 13]
Data in table (5) confirms that there were significant interaction between geno-
types and cutting treatments on all yields and yield components characters with the
exception of spick yield, No. of Spike m2 , 1000- grain weight and average spick
weight. The highest number was exhibited by the combination of genotype (2) and
non cutting treatment in most characters.
Dry forage yield significantly affected frame vary genotypes. Data in Table 6
shows highly significant differences between genotypes on dry forage yield; the
maximum dry forage yield value recorded by genotype (2) is 2616.33 Kg ha-1,

45
 Table 3: The means of genotypes for yield and yield components
Spike No. Grain 1000- Average
Biological Straw
Genotypes yield of yield grain spike
yield yield HI
(G) −1 (kg Spike (kg −1 weight weight
(kg ha ) (kg ha )
ha−1 ) (m−2 ) ha−1 ) (g) (g)
1 15777.73 7731.15 485.83 5750.28 10027.55 0.36 37.75 1.887
2 15277.80 8006.90 512.67 6405.55 8872.12 0.42 40.10 2.02
3 12577.60 6678.95 537.33 5111.78 7465.82 0.41 39.46 1.51
4 12699.93 6050.92 426.17 4636.52 8063.42 0.37 35.98 2.15
5 9155.47 4045.67 310.50 3154.77 5990.82 0.35 33.31 1.89
6 13850.93 6040.83 460.33 4922.82 8928.12 0.36 41.13 2.47
LSD
(P ≤ 3834.84 1679.76 130.61 1482.06 n.s n.s 1.68 0.37
0.05)

Table 4: The effect of Cutting treatment on yield and yield component

Spike No. Grain 1000- Average
Biological Straw
Cutting yield of yield grain spike
yield yield HI
treatments (kg ha−1 ) (kg Spike (kg −1 weight weight
(kg ha )
ha−1 ) (m−2 ) ha−1 ) (g) (g)
N.CU 14918.47 7255.18 473.61 5747.24 9171.18 0.39 38.80 2.06
CU. 11528.02 5596.29 437.33 4246.67 7278.09 0.37 37.11 1.91
LSD
(0.05) for 373.26 221.96 21.01 152.75 412.05 0.02 0.86 n.s
CU.

Table 5: The effect of interaction between genotypes and cutting treatments on yield and yield
component
Biological Spike No. Grain Straw 1000- Average
Genotypes Cutting yield yield of yield yield grain spike
HI
(G) treatments (kg (kg Spike (kg (kg weight weight
ha−1 ) ha−1 ) (m−2 ) ha−1 ) ha−1 ) (g) (g)
N.CU 18644.40 8642.83 488.33 6682.47 11961.93 0.353 39.03 2.05
1
CU. 12911.07 6819.47 483.33 4818.10 8093.17 0.373 36.46 1.72
N.CU 18177.73 9093.17 525.67 7252.43 10925.03 0.400 41.03 2.04
2
CU. 12377.87 6920.63 499.67 5558.67 6819.20 0.447 39.16 2.00
N.CU 14333.07 7630.30 572.00 6030.03 8303.03 0.427 40.23 1.68
3
CU. 10822.13 5727.60 502.67 4193.53 6628.60 0.393 38.70 1.34
N.CU 13733.33 6709.50 443.00 5254.20 8479.13 0.390 36.67 2.07
4
CU. 11666.53 5392.33 409.33 4018.83 7647.70 0.350 35.29 2.22
N.CU 9800.00 4760.70 330.00 3760.30 6039.70 0.387 34.40 1.84
5
CU. 8510.93 3330.63 291.00 2549.23 5941.93 0.310 32.22 1.94
N.CU 14822.27 6694.57 482.67 5504.00 9318.27 0.377 41.42 2.69
6
CU. 12879.60 5387.10 438.00 4341.63 8537.97 0.347 40.83 2.24
LSD (P≤0.05) 914.285 n.s n.s 374.16 1009.30 0.04 n.s n.s

46
 Table 6: The effect of genotypes on dry forage yield
Dry forage yield
Genotypes (G)
(kg ha−1 )
1 2051.67
2 2616.33
3 1460.67
4 1108.14
5 1604.61
6 1939.41
LSD (P≤0.05) 518.53

while the minimum dry forage yield value recorded by the genotype (4) is 1108.14
Kg ha-1.Simlar results were reported previously which confirmed that wheat yield
and its components affected were due to the differences between varieties (18, 19,
20 and 21).

References

[1] Ullrich, S. E. 2011. Barley: Production, Improvement, and Uses. Press Black-
well Publishing Ltd.
[2] Grando, S. and M.H. Gomez. 2005. Preface. In S. Grando & H. Gomez
Macpherson (Eds.), Food barley: Importance, uses and local knowledge (pp. ixŰx).
Syria: ICARDA
[3] Badr, A., Müller, K., Sch?fer, R., El Rabey, H., Ibrahim, H.H., Pozzi, C., Ro-
hde, W., and Salamini, F. 2000. On the origin and domestication history of barley
(Hordeum vulgare). Mol. Biol.
[4] Dumphy, D. S. M. E and Hoit, E. C. 1982. effect of forage utilization on wheat
grain yield. Crop sci. 22: 160-109.
[5] Mohammed- Mahfoth. A. 1990. Effect of different stages of tipping on
yield and its component for two varieties of barley under condition of north Iraq.
Mesopotamia. J. of Agri. 22. 1: 331- 349.
[6] Royo. A. C. and f. Tribb. 1997. Triticale and barley for grain and for dual-
purpose (forage and grain) in a Mediterranean-type environment 11. yield and yield
component and quality. Aust. J. Agric. Res. 48: 423-32.
[7] Hadjicnristodoulou. A. 1986. Brley improvement for different uses in dry area
in cyrus. Rachis 8. 1: 9-11.
[8] Droshiotis, D. N. and L. Wilman. 1987. effect of harvesting program sowing
date on forage yield, digestibility nitrogen concentration, tillers area fraction barley
incyprus. J. Agri. Sci. comb.109: 95-105.
[9] Al- Hasan. Abbas. M. A. 1995. Effect of growth stage, clipping and seed-
ing rateson growth, forage and seed production of barley under rain-fed conditions.
Ph.D. thesis submitted to college of Agriculture and forestly.Mosul University.
[10] Morrill. L. Ag. 1973. Effect of fertilization and other management variable
on wheat forage Res. Rep. No. p: 684.
[11] Yau. S. K. 2003. Yield of early planted barley after clipping or Grazing in a
semiarid area. Agro. J. Vol. as 4: 821-827.

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 [12] Hasan. Abbas. M. A. Mohammed. 2008. Forage and grain yield of same bar-
ley genotypes under rain-fed conditions. Mesopotamia-s. J. of Agri. 361: 161-166.
[13] JMP, Version 7. ŞSAS Institute Inc.Ť, Cary, NC, (1989-2007).
[14] Al-Rawi- Baha. A. and A. M. Shamma. 1991. Effect of processes clipping on
morphological yield components and grain yield of barley. (H Űvulgar L).J. of IPA
Agricultural Research Center. Vol. (1). 1: 16-31.
[15] Al-Shamma. A. M. B. A., AlRwi and M. J Abdulrahman. 1993. Nature of
tolerance to continuous clipping among three Iraqi fodder barley land races and
some of their extracts. Proceeding of the workshop on technology transfer in the
production of cereals and legumes 20-22 Sep. 24-38.
[16] Yesmin,SH., Moushumi , A., and Belal, H. Yield and Seed quality of barley
(Hordeum Vulgare) as affected by variety, nitrogen level and harvesting time.In. J.
Agric. Sci , Vol.7, No.5 : 262-268. 2014.
[17] El -Banna, Mohamed Naguib .2011.Evaluation of 16 barley genotypes under
calcareous soil conditions in Egypt. J. Agric. Sci , Vol.3, No.1 : 105-121.
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wheat yield in the mid South .J.Prod.Agric.3:345-348.
[19] Musaddique. M., A. Hussain, S. Aftab wajid and A. Ahmadm 2000.Growth,
Yieldand Components of yield of different Genotypes of Wheat. Deparment of
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48
 Foliar Application of Humic Acid, Iron and Sprays Number on Chemical
Quality of Grape (Vitis vinifera L.) cv. TAIFI∗
Shawkat Mustafa Al-Atrushy1 and Shayan Akram Mustafa2
1
Hort. Department, College of Agriculture, Duhok University, Kurdistan-Iraq
2
Hort. Department, Faculty of AAgriculture Sciences, Sulaimani University,
46001 Sulaimani, Kurdistan-Iraq
Abstract
This study was carried out on 12-years old Grapevine (Vitis vinifera L.) cv.
Taifi grown in Bara-Buhar village- near Zawita town, Duhok governorate- Kurdis-
tan region/Iraq, during the growing season 2012, in order to investigate the effect of
number of sprays, foliar application of Humic acid at three concentrations (control,
2 and 4 g.L− 1) and three concentrations of iron (control, 100 and 200 mg.L− 1) as
well as their interactions on chemical characteristics of grape. The foliar applica-
tion was done once at two weeks before full bloom and twice carried out at berries
swelling. The results showed that sprays twice, foliar application of Humic acid
especially at 4 g.L− 1 and Iron especially at concentration 200 mg.L− 1 significantly
increased TSS, total sugar % and juice %, whereas only application of Humic acid
and Iron significantly increased anthocyanin content. All treatment significantly re-
duced total acidity percentage. The interaction between number of spray, Humic
acid and iron significantly enhanced all detected traits, since vine spraying twice +
4 g.L− 1 Humic acid + 200 mg.L− 1 Iron significantly outshine in chemical charac-
teristics of grape compared to the untreated trees (control).
Keywords:spray number, Humic, iron, quality, grape.

Introduction

Grape (Vitis vinifera L.) belongs to Vitaceae family, is perhaps the most widely
cultivated fruit crop of the world in varying climatic zones extending from the
temperate to the tropics region (Yogeesha, 2005). According to the Food and Agri-
culture Organization (FAO, 2012) 7,586,600 square km of the world is dedicated
to grapes, producing over 8519148 tons of grapes. Approximately 71% of world
grape production is used for wine, 27% as fresh fruit, and 2% as dried fruit. A
portion of grape production goes to producing grape juice to be reconstituted for
fruits canned ”with no added sugar” and ”100% natural”. The area dedicated to
vineyards is increasing by about 2% per year.
Humic acid is a commercial product contains many elements which increas-
ing the availability of nutrient elements and consequently affects plant growth and
yield. Humic acid reduces other fertilizer requirements, increase the yield, pro-
tein and mineral contents of most crops and establish a desirable environment for
microorganism development (Ferrara and Bunetti, 2010).
In plant, micronutrients are playing an important role in the production and pro-
ductivity. Among micronutrients, iron plays vital role in synthesis of chlorophyll,
carbohydrate production and cell respiration, chemical reduction of nitrate and sul-
fate and in nitrogen assimilation (Mengel and Kirkby, 2000 and Havlin, et al., 2005).

49
 Therefore, the present study was undertaken to investigate the effect of number of
sprays, Humic acid and iron on yield and quality parameters of grape cv. Taifi.
Numerous researchers studied the effect of various concentrations of Humic acid
and Iron on grape in aim to improve its yield and quality. [Al-Atrushy (2012)
investigated the effect of foliar application of three concentrations of Humic acid
(0, 1 and 3 g.L−1 ) and three concentration of Bread yeast (0, 2 and 4 g.L−1 ) on
grapevine cv. Rash-Mew. Results indicated that Humic acid at 2 and 3 g.L-1
was accompanied with improving berries characteristics which increased TSS, total
sugar and juice density compared to the control, also the same traits also lead to
decrease in total acidity percentage in berry juice. Shaheen, et al. (2012) studied
the effect of two different compost types, plant residues and plant + animal residues
with two natural rocks, rock phosphate and feldspar respectively, all used treatments
were applied with or without NPK bio-fertilizers and humic acid. Result indicated
that the compost plus bio-fertilizer and humic acid fertilizers improved fruit quality
expressed by increasing TSS, total sugars content and decreasing acidity. Therefore,
these organic and neutral fertilizers in combination with NPK, bio-fertilizers and
humic acid can be recommended for Crimson grapevine to improve productivity
and quality and produce a healthy product. Rui, et al. (2012) studied the effect of
some fertilization with (Ca, Mg, S and Humic acid) on chemical characteristics of
grape. Results showed that the Humic acid significantly increased the content of
soluble solid and total sugar.
Ismail, et al. (2010) used some different treatments with magnetic iron ore, metal
compound fertilizers and bio-NK as a biocide were applied as soil amendments to
study their effects on chemical properties of the grapevine berries of two grapevine
cultivars Superior and Thompson Seedless. Results showed that the chemical prop-
erties of the two cultivars were significantly affected by the treatments, total soluble
solids (TSS%) was increased by all treatments compared to control, all treatments
gave low total acidity (TA) in juice and the higher value (0.89%) was in the un-
treated plants. Ahmed, et al. (2011) investigated the effect of some amino acids,
nutrients Mg, Zn, Fe, P, K bio-fertilizers and EM treatments on yield and quality of
the berries of Thompson Seedless Grapevines. Results indicated to significant in-
crease in total soluble solids and total sugars and decrease in total acidity compared
to the control.

Material and Methods

The present study was conducted during growing season of 2012 on a private vine-
yard located at Bara-Buhar village, Duhok governorate, Kurdistan region, Iraq, in
order to investigate the effect of sprays number of three concentrations of both Hu-
mic acid (0, 2 and 4 g.L−1 ) and Iron as FeSO4 .7 H2 O (0, 100 and 200 mg.L−1 ), on
the chemical characteristics of grape cv. Taifi, irrigated with drip irrigation system.
Therefore, the experiment was consisted of 18 interaction treatments (two dates
of spraying × 3 conc. of Humic acid × 3 conc. of Iron). A randomized complete
block design with three factors was followed in the experiment. Every treatment
consisted of one vine per replicate with three replications, so the numbers of vines

50
used were 54 vines. The investigated vines were trained on (T) trellis, vines were
pruned by leaving 7 fruiting canes each with 8 eyes and four renewal spurs each with
two eyes (the total eyes left on each vine were 64 eyes). The experimental vines
were chosen to be uniform as possible in vigor (measurement of trunk diameter at
30 cm above the soil) and similar in size.
These vines were planted at 2.5 × 2.5 m inter and intra rows, spaces and had
12 years old. All vine undertaken in this study received the regular agricultural
and horticultural practices that usually carried out in the commercial vineyard. The
harvest date had been recorded after the TSS of berries reached 16% and the color
changed to rose or purple.
Potential effects of sprays number, Humic acid and iron were evaluated in terms
of the change in TSS, total sugar % and juice %, total acidity percentage and
anthocyanin content.
All results were analyzed using SAS program (SAS 2003), the means were com-
pared by DuncanŠs multiple range tests at 5% level.

Results and Discussion

1. Total soluble solids percentage (TSS %)

Table (1) illustrates that Increasing number of sprays significantly increased TSS%,
the highest TSS (16.44 %) is registered in vine sprayed twice.
The same table shows data of vines sprayed with different Humic acid concen-
trations. The results showed that (4 g.L−1 ) gave the highest significant value of TSS
(17.14 %) which was significantly higher than that of 2 mg.L−1 (16.00%) and both
of them were significantly higher than control (15.51%).
Effect of foliar application of Iron on TSS % of grape berries at harvest is
shown in table 1, which indicates that spraying Fe at conce. (200 mg.L−1 ) recorded
(17.04 %) which was significantly superior to berries TSS (100 mg.L−1 ) which also
significantly differed from (0 mg.L−1 ) which recorded (15.22 %).
For the interaction, the results show that one spray and (4 g.L−1 ) gave the higher
significant TSS (17.22 %), the interaction of number of sprays and foliar application
of Iron had significant effect on TSS %, the highest value (17.11 %) was obtained
from the interaction of twice sprays and (200 mg.L−1 ) compared with the lowest
TSS (14.97 %) from control.
Regarding the interaction of spraying HA and Fe, data in table1 shows that the
interaction of 4 g.L−1 HA + 200 mg. L−1 Fe recorded the highest TSS (17.88 %),
which was differed significantly from all the other interactions.
The interactions of the three studied factors, treatment represented that combi-
nation of control was significantly lower in TSS (13.67 %) than combination of one
spray+4 g.L−1 HA+ 200 mg.L−1 Fe which recorded (18.08%). Means within a col-
umn, row and their interactions followed with the same letters are not significantly

51
 Table 1: Effect of No. of sprays, Humic acid, Iron and their interaction on total soluble solid (TSS
%) of grape cv. Taifi.
Number of Humic Acid Iron (g.L−1 ) Sprays Main Effect
−1
Sprays (g.L ) X HA of sprays
0 100 200
0 13.67 f 15.50 de 16.17 cd 15.11 c
Once 2 14.67 ef 15.67 c-e 16.67 b-d 15.67 bc 16.00 b
4 16.58 b-d 17.00 a-c 18.08 a 17.22 a
0 14.42 ef 16.67 b-d 16.67 b-d 15.92 b
Twice 2 15.50 de 16.50 b-d 17.00 a-c 16.33 b 16.44 a
4 16.50 b-d 17.00 a-c 17.67 ab 17.06 a
Sprays Once 14.97 c 16.06 b 16.97 a Main Effect
X Iron Twice 15.47 bc 16.72 a 17.11 a of HA
0 14.04 e 16.08 c 16.42 bc 15.51 c
HA X
2 15.08 d 16.08 c 16.83 bc 16.00 b
Iron
4 16.54 bc 17.00 b 17.88 a 17.14 a
Main effect of Iron 15.22 c 16.39 b 17.04 a

different from each other according to Duncan multiple ranges test at significant
level 5%.

Total sugar percentage (%)

Table 2 clearly shows that increasing number of sprays of Humic acid or Iron,
significantly increased total sugars %, the highest total sugar % (16.70 %) was
obtained from the vine receiving twice sprays of both factors as compared with one
spray of both factors (13.68 %).
The data in table 2 displays, the highest value of total sugar % was obtained with
spraying Humic acid (4 g.L−1 ) which was (16.61 %) and Fe (200 mg.L−1 ) which
was (16.63 %).
Pertaining to the interaction of number of sprays + HA, number of sprays + Fe
and Humic acid + Iron, data in table 2 shows that the interaction of twice sprays
+ (4 g.L−1 ), twice sprays + 200 mg.L−1 Fe and 4 g.L−1 HA + 200 mg.L−1 Fe
gave the higher total sugar % (18.46, 18.32 and 18.35 %) respectively. Means
within a column, row and their interactions followed with the same letters are not
significantly different from each other according to Duncan multiple ranges test at
significant level of 5%.
The same table also shows that the interaction of twice sprays + 4 g.L−1 HA +
200 mg.L−1 Fe (20.14 %) was distinctive significantly from most of the other three
studied factors interactions.

Total acidity percentage (%)

Table 3 shows that there are significant differences between the treatments. In-
creasing number of sprays of HA or Fe significantly decreased the total acidity
percentage, the lowest acidity (0.45 %) was obtained in the vine sprayed twice as
compared to the highest acidity (0.50 %) with vine sprayed once.

52
Table 2: Effect of No. of sprays, Humic acid, Iron and their interaction on total sugar percentage
(%) of grape cv. Taifi.
Number of Humic Acid Iron (g.L−1 ) Sprays Main Effect
−1
Sprays (g.L ) 0 100 200 X HA of sprays
0 11.64 f 12.44 ef 13.16 d-f 12.41 d
Once 2 13.21 d-f 13.30 d-f 15.07 c-f 13.86 cd 13.68 b
4 13.39 d-f 14.30 c-f 16.57 b-d 14.75 bc
0 12.02 ef 17.20 a-c 17.55 a-c 15.59 bc
Twice 2 15.19 b-f 15.70 b-e 17.27 a-c 16.05 b 16.70 a
4 16.50 b-d 18.76 ab 20.14 a 18.46 a
Sprays Once 12.75 c 13.35 bc 14.93 b Main Effect
X Iron Twice 14.57 bc 17.22 a 18.32 a of HA
0 11.83 c 14.82 b 15.35 b 14.00 b
HA X
2 14.20 b 14.50 b 16.17 ab 14.96 b
Iron
4 14.95 b 16.53 ab 18.35 a 16.61 a
Main effect of Iron 13.66 c 15.28 b 16.63 a

Results show highly significant differences in acidity % due to spraying different

concentrations of HA solution. Increasing HA concentration from 0 to 4 g.L−1
significantly decreased total acidity percentage, while vines sprayed with HA at
concentration 2 g.L−1 had no significant effect on total acidity percentage in Berry
juice.
The same table illustrates that spraying vines with Iron had no significant on
acidity.
It’s clear from table 3 that the acidity % values of grape juice was affected by the
combination of number of sprays + HA, sprays + Fe and HA + Fe, The minimums
value of total acidity (0.40, 0.43 and 0.44 %), was recorded from vine sprayed
twice + 4 g.L−1 HA, twice + 200 mg.L−1 Fe and 4 g.L−1 HA + 200 mg.L−1 Fe
respectively. Means within a column, row and their interactions followed with the

Table 3: Effect of No. of sprays, Humic acid, Iron and their interaction on total acidity percentage
(%) of grape cv. Taifi.
Number of Humic Acid Iron (g.L−1 ) Sprays Main Effect
−1
Sprays (g.L ) 0 100 200 X HA of sprays
0 0.60 a 0.58 ab 0.50 a-d 0.56 a
Once 2 0.50 a-d 0.44 cd 0.45 b-d 0.46 b-d 0.50 a
4 0.49 a-d 0.46 b-d 0.50 a-d 0.48 bc
0 0.54 a-c 0.50 a-d 0.49 a-d 0.51 ab
Twice 2 0.44 cd 0.44 cd 0.41 cd 0.43 cd 0.45 b
4 0.42 cd 0.40 cd 0.39 d 0.40 d
Sprays Once 0.53 a 0.49 ab 0.48 ab Main Effect
X Iron Twice 0.47 ab 0.44 b 0.43 b of HA
0 0.57 a 0.54 ab 0.49 a-c 0.53 a
HA X
2 0.47 bc 0.44 c 0.43 c 0.45 b
Iron
4 0.45 bc 0.43 c 0.44 c 0.44 b
Main effect of Iron 0.50 a 0.47 a 0.46 a

same letters are not significantly different from each other according to Duncan
multiple ranges test at significant level of 5%.

53
 Table 3 shows the interaction effects of three factors on total acidity, the lowest
significant value of total acidity % of grape juice (0.39 %) was recorded from the
combination of (twice sprays + 4 g.L−1 HA + 200 mg.L−1 Fe) compared to the non
sprayed one which recorded (0.60 %).

Juice percentage (%)

Data in table 4 shows that increasing in number of sprays of humic acid or iron
significantly increased juice percentage (%) of grapevine cv. Taifi. The highest
juice percentage (83.51 %) was obtained from spraying twice of both humic acid
and iron compared with the lowest value of the juice percentage (72.40 %) with
spraying once of both factors.

The highest significant juice % (79.58 and 80.02 %) resulted from grapevine
treated with 4 g.L−1 HA and 200 mg.L−1 Fe respectively, which was significantly
higher than control.

Table 4 clearly shows the interactions between sprays + Humic acid, sprays and
Fe and Humic acid + Fe on increasing the juice percentage. The highest value
(85.29, 84.74 and 81.20%) was obtained from the interaction of twice sprays + 4
g.L−1 HA, twice sprays + 200 mg.L−1 Fe and 2 g.L−1 Humic acid and 200 mg.L−1
Iron respectively. Means within a column, row and their interactions followed with

Table 4: Effect of No. of sprays, Humic acid, Iron and their interaction on juice percentage (%) of
grape cv. Taifi.
Number of Humic Acid Iron (g.L−1 ) Sprays Main Effect
−1
Sprays (g.L ) 0 100 200 X HA of sprays
0 63.17 f 69.15 e 73.63 de 68.65 c
Once 2 71.30 e 75.01 c-e 77.77 b-d 74.69 b 72.40 b
4 72.85 de 74.24 de 74.50 de 73.86 b
0 82.32 ab 83.02 ab 82.73 ab 82.69 a
Twice 2 80.77 a-c 82.24 ab 84.62 a 82.54 a 83.51 b
4 84.72 a 84.29 a 86.87 a 85.29 a
Sprays Once 69.11 c 72.80 b 75.30 b Main Effect
X Iron Twice 82.60 a 83.18 a 84.74 a of HA
0 72.75 c 76.08 bc 78.18 ab 75.67 b
HA X
2 76.04 bc 78.63 ab 81.20 a 78.62 a
Iron
4 78.78 ab 79.27 ab 80.68 a 79.58 a
Main effect of Iron 75.85 b 77.99 ab 80.02 a

the same letters are not significantly different from each other according to Duncan
multiple ranges test at significant level of 5%.

Combination between the three studied factors affected significantly on the juice
% as shown in table 4. The highest value (86.87 %) was recorded from combination
of twice sprays + 4 g.L−1 HA + 200 mg.L−1 Fe, which differed significantly from
lowest value (63.17 %) of the control (one spray + 0 g.L−1 HA + 0 mg.L−1 Fe).

54
Anthocyanin (mg/100g fresh weight)

The effect of number of sprays of Humic acid or Iron showed no significant effect
on the Anthocyanin content in the juice of grapevine cv. Taifi.
The data in table 5 revealed that spraying grapevine with Humic acid concentra-
tions had a significant effect on Anthocyanin, especially at (2 g.L−1 ), which gave
the maximum value (0.24 mg/100g fresh wt.) while (0 g.L−1 ) gave the minimum
value (0.20 mg/100g fresh wt.).
Increasing the concentration of the Iron increased the Anthocyanin the highest
value (0.25 mg/100g fresh wt.) was registered in the vine foliar sprayed with (200
mg.L−1 ) Iron, which surpasses significantly on the other treatment.
The interaction of twice sprays + 2, 4 g.L−1 HA, twice sprays + 100, 200 mg.L−1
Fe and 2 g.L−1 HA + 200 gm.L−1 Fe were recorded significantly maximum An-
thocyanin content, which was (0.24, 0.25 and 0.26 mg/100g fresh wt.) respectively.
Means within a column, row and their interactions followed with the same letters

Table 5: Effect of No. of sprays, Humic acid, Iron and their interaction on Anthocyanin (mg/100g
F.Wt.) of grape cv. Taifi.
Number of Humic Acid Iron (g.L−1 ) Sprays Main Effect
−1
Sprays (g.L ) 0 100 200 X HA of sprays
0 0.11 d 0.15 cd 0.27 a 0.18 b
Once 2 0.27 a 0.21 a-d 0.25 ab 0.24 a 0.21 a
4 0.22 a-c 0.21 a-d 0.21 a-d 0.21 ab
0 0.15 b-d 0.26 a 0.23 a-c 0.21 ab
Twice 2 0.21 a-d 0.24 a-c 0.27 a 0.24 a 0.23 a
4 0.22 a-c 0.24 a-c 0.26 a 0.24 a
Sprays Once 0.20 a-c 0.19 c 0.24 ab Main Effect
X Iron Twice 0.19 bc 0.25 a 0.25 a of HA
0 0.13 b 0.21 a 0.25 a 0.20 b
HA X
2 0.24 a 0.22 a 0.26 a 0.24 a
Iron
4 0.22 a 0.23 a 0.23 a 0.23 ab
Main effect of Iron 0.20 b 0.22 ab 0.25 a

are not significantly different from each other according to Duncan multiple ranges
test at level of significant 5%.
For the interaction data in table 5 indicated that the interaction of Number of
sprays + Humic acid + Iron were differed significantly, the highest Anthocyanin
value (0.27 mg/100g fresh wt.) was obtained from the interaction of twice sprays
+ (2 g.L−1 ) HA + (200gm.L−1 ) Fe, which was significantly differed from many of
the interaction between studied factors
It’s cleared from (Tables 1, 2, 3, 4 and 5) that twice sprays significantly in-
creased most of the chemical characteristics of grape cv. Taifi (Total soluble solid,
Total sugar percentage, juice percentage and Anthocyanin) whereas, the mentioned
treatments significantly reduced Total acidity percentage.
The reason of this effect may be attributed to the role of Humic acid and Iron in
activation of vegetative growth characteristics of grapevine which in turn reflected

55
 on the supplement of more accumulation of carbohydrates in plants and resulting
in the stimulation of ripening (Zhu and Zhu, 2000).
Ferrara and Brunetti (2010) state that Humic acid application at various steps
have significant effect in improving quality parameter of berries such as reducing
total acidity, increasing TSS and total sugar percentage, so they confirmed that
Humic acid, if applied twice can increase chemical quality of table grape and can
find appositive application in an organic and sustainable viticulture also the role
of Humic acid in achieving a good balance between growth and fruiting which is
reflected in turning on increasing accumulation of carbohydrate and thus stimulation
of repining (Li-Nan, et al.,1999 and Zhu and Zhu, 2000).
Yogeshaa (2005) mentioned that twice spray of Iron helped in continuous supply
of Iron to the plant which reflected in enhancing metabolic process in plant. Also,
its role in increasing of photosynthesis product led to increase the accumulation
of carbohydrates in Berries then improves chemical characteristic of Berries (Al-
Eshaqi, 2007).

References

Ahmed, F. F.; A. A. Ibrahiem; A. E. M. Mansour; E. A. Shaaban and M.S. El-

Shammaa. 2011. Response of Thompson Seedless Grapevines to Application
of Some Amino Acids Enriched with Nutrients as Well as Organic and Biofer-
tilization Research J. Agric. and Biological Sci. 7(2): 282-286.
Al-Atrushy, Sh. M. M. 2012. Effect of foliar application of Humic acid and Bread
yeast on some vegetative growth characteristics, yield and quality of grape
(Vitis vinifera L.) cv. Rash-Mew under non irrigation condition. J.of Duhok
Univ. 15(1).
Al-Eshaqi, J. M. Kh. 2007. Effect of NPK fertilization and foliar application of
iron, gibberelic acid on the growth and yield of Salamy pomegranate cultivar.
Ph. D. Thesis, College of Agric. Univ. of Mosul, Iraq.
FAW Food and Agricultural Organization. 2012. Production year book, of United
Nation, Rome. Italy.
Ferrara, G. and G. Brunetti. 2010. Effects of the times of application of a soil
humic acid on berry quality of table grape (Vitis vinifera L.) cv Italia. Spanish
J. of Agric. Res. 8(3): 817-822.
Havlin, L. L.; J. D. Beaton; S. L. Tisdale and W. L. Nelson. 2005. Soil Fertility
and Fertilizers. 7th ed. Upper Saddle River, New Jersey 07458.
Ismail, A. E.; S. S. Soliman; E. M. Abd El-Moniem; M. S. Awaad and A. A.
Rashad. 2010. Effect of magnetic iron ore, metal compound fertilizer and
bio-NK in controlling root-knot nematode of grapevine a newly reclaimed area
of Egypt. Pak. J. Nematol. 28(2): 307-328.
Li-Nan, X.; B. L. Wang and B. W. Lu. 1999. Study of the effect of apple liquid
fertilizer on the growth and fruit development of Starkrinson apple variety.

56
 China- Fruits. 4: 20-21.
Mengel, K. and E. A. Kirkby. 2000. Principles of plant nutrition. Trans. Withs.
N. Al-Niem. Dar Al-Kutud for Printing and Publishing. Mosul. Univ. Iraq.
(In Arabic).
Rui, W.; S. Quan; Z. Xiao-juan; G. Jie; S. Hai-li and Ch. Qing-rui. 2012. Influence
of reasonable irrigation and fertilization on early fruit wine grape in Helan
Mountain. Water saving irrigation. 06. (Abstract).
SAS Institute. 2003. Statistical Analysis System Procedures Guice, Version 9, third
ed. Institute Inc. Cary.
Shaheen, M. A. G.; S. M. Abdel-Wahab; E. A. Hassan and A. M. R. A. AbdelAziz.
2012. Effect of some soil conditioners and organic fertilizers on vegetative
growth and quality of Crimson seedless grapevines. J. of Hortic. Sc. &
Ornamental Plants 4(3): 260-266
Yogeesha. 2005. Effect of iron on yield and quality of grape (Vitis vinifera L.) in
calcareous vertisol. M.Sc. Thesis, Univ. of Agric. Sci. Dharwad.
Zhu, R. F. and R. E. Zhu. 2000. Effect of (Opal) HA-2002 liquid fertilizers on the
growth and fruiting of pears trees. South China Fruits. 29(2): 43- 47.

57
58
 Induction of Sister Chromatid Exchange in Lymphocytes of Brick Kiln
Workers in Baghdad
Zainab M.T.Jaafar , Nazar A.Auda2 and Sa’adyia O. Mohmmed2
1
1
Directorate of Agriculture Research /Ministry of Science and Technology.
Baghdad, Iraq P.O.Box765
2
Department of Biology, College of Science, Baghdad University, Baghdad, Iraq
[email protected]
Abstract
This study was conducted on the brick kilns at Al-Nahrawan area south of Bagh-
dad city. To investigate the effect of the pollutant emission (fumes and gases )
released from burned fuel in the kiln, blood samples were taken from group indi-
viduals of workers at this brick kilns and non working individuals as control. The
influence of these emissions on the workers health was measured by examining the
sister chromatid exchange (SCEs) frequency as cytogenetic parameters. The data
showed significant differences between workers in brick kiln and healthy individ-
uals (9.96,4.38) respectively, the relation ship of age groups and SCEs frequency
demonstrate that 30-39 years and over gave higher value (11.4) and was significant
with age groups of 10-19 and 20-29 years. the data showed relationship between
smoking workers and non smoking which gave high value of SCEs frequencies
(14.34) compared with non smoking and control (9.15, 4.34) respectively.
Keywords: Sister chromatid exchange. Human blood lymphocytes, brick kiln work-
ers.

Introduction

Brick manufacturing is a labor intensive informal industry of India employing chil-

dren as major workforce. This unorganized sector is growing multifold due to ever
increasing demands of brick for continued urbanization and rural expansions with
techno-modern concretization. Despite of engagement of large population in brick
industries data on number of employed workers, accidents, health and safety issues
are missing. Brick making is a painstaking labor job often performed by families
accompanied by children living in temporary settlements ( Rabin , 2015). The
brick kilns emit toxic fumes containing suspended particulate matters rich in carbon
particles and high concentrations of carbon monoxides and sulfur oxides that are
harmful to eye, lungs, throat and stunt the mental and physical growth of children
(Saber et al ., 2014 ). The intensity and size of the pollutants emitted by the brick
kilns are varied according to several factors include: size and age of the brick kilns,
the system of maintenance, size of production and availability of pollution control
system. Most of these factories are committed in employing a mechanical burning
technology and depend mainly on primitive methods of burning the fuel which lead
to incomplete burning which result in spreading pollutants gases and fumes emitted
by the brick kilns fuel burning released through the chimney and cracks in the walls
(Abidand Salman,2005). Sister chromatid exchange is the process whereby, during
DNA replication, two chromatids break and rejoin with one another, SCEs technique
is considered a very important test to check potential damage to the DNA resulting

59
 from contact with a mutagen. This test can be used not only for estimation of influ-
ence of genotoxic factors on chromosome, and to prevent and monitor diseases but
also can have more global applications, allowing for assessment of such factors as
sex, age, and race, as well as for assessments of genetic immunity. (Szeleszczukl, et
al .,2014). It enables the detection of early genotoxic effects at the level of DNA and
chromosome,(Wojeik et al., 2012). The process is considered to be conservative and
error Űfree, since no information is generally altered during reciprocal interchange
by homologous recombination. Today, cytogenetics provides several biomarkers for
chromosomal instability assessment, one of which is the sister chromatid exchange
(SCE) frequency in cells. SCE is a natural process that implicates the exchange of
homologous genetic segments as a mode of repair mechanism. The methodology
of SCEs has been proved to be a very useful tool with predictive value, for de-
tecting harmful effects on DNA, caused by various physical and chemical factors.
Increased frequency of this index reflects the existence of genotoxicity in cells and
the subsequent failure of repair mechanisms to recover the damaged site. Several
studies have reported that SCE analysis is a very sensitive method, able to detect
mutagens and/or carcinogens( Wilson and Thompson 2007). Chromosome ends
are known hotspots of meiotic recombination and double-strand breaks. monitored
mitotic sister chromatid exchange (SCE) in telomeres and subtelomeres and found
that 17% of all SCE occurs in the terminal 0.1% of the chromosome. Telomeres
and subtelomeres are significantly enriched for SCEs, exhibiting rates of SCE per
basepair that are at least 1,600 and 160 times greater, respectively, than elsewhere
in the genome( Rudd, et al.,2007).

The aims of study is to investigate the pollutants emitted by brick kilns through
burning fuel and examine their effects on the cellular disorders of the workers by
using SCEs as cytogenetic marker .

Materials and Methods

Data collection

It was done from brick kiln workers at Al-Nahrawan area at Baghdad city which
they exposed to pollutants of brick kiln products.

Blood samples

From each subject, (4ml) of peripheral blood were collected under aseptic conditions
by vein- puncture using disposable syringe promoted with heparin;. The blood was
placed in a cool- box and transferred to the laboratory. The samples were used for
cytogenetic assay of sister chromatid exchange (SCEs).

60
Cytogenetic analysis of human blood lymphocytes

Blood cell culture for cytogenetic analysis was performed following the procedure:
Human peripheral Blood was collected into heparin coated syring by taking 0.25ml
in test tube containg 2ml of RPMI -1640. Supplemented with 10% fetal calf serum
and antibiotic solution of penicillin and streptomycin With 0.25ml of PHA and
10ug/ml of 5BrDU mix gently, then incubated at 37◦ C four 72 hrs. Adding of 0,2
ml of colcemid before 2hrs of harvesting.The test tubes were centrifuged at speed of
(2000 rpm) for10 min.The supernatant was removed and (5 ml) of KCl (hypotonic
solution 0.075 M) was added, and then the test tubes were kept in the incubator
at (37◦ C) for one hour with shaking.. The tubes were centrifuged at (2000 rpm)
for 10 min.The supernatant was removed and the Carney’s fixative was added as
drops on the inside wall of the tubes with continuous shaking. the cells were kept
at (4◦ C) for 30 min for fixation. The tubes were re-centrifuged at (2000 rpm) for
10 min.. The process was repeated 3 times and after that, the cells were suspended
in 2 (ml) of fixative solution..A few drops from the tubes were dropped vertically
by the use of Pasteur pipette from the height of 3 feet at a rate of (4 -5) drops to
give the chance for chromosomes to be spread well, then the slides were kept at
room temperature.The slides were stained with DAPI stain for 10 min and then take
the slides from the stain and add 2-3 drops of DAPI buffer and put the cover slide
and examine by fluorescent microscope to count the number of SCEs.(Jafar et al.,
2001).

Statistical analysis

One way analysis of variance was performed to test whether group variance was
significance or not, the comparison between groups were used analysis of variance
(Al- Mohammed et al ., 1986).

Results and Discussions

The results of SCEs induction on the blood lymphocytes in workers of the brick
kiln

SEC is known to result from reciprocal DNA interchange in homologous loci

at certain rates in all cells. The results of SCE presented in table . showed the
comparison between brick kiln workers and control healthy individuals in induction
of SCE frequency per cell, it noticed that average of SCE frequency increased in the
group of brick kiln workers, it gave 9.69±1.06 while in group of healthy individuals
gave 4.34±0.15 SCE /cell . The statistical analysis by using t-test showed significant
at probability P (≤ 0.01).

61
 Table 1: Compare between brick kiln workers and the control in SCEs induction
Group SCE’s ± SD
Workers 9.96 ± 1.06
Control 4.38 ± 0.15
T-test value 2.219 **
P-value 0.0038
** (P<0.01)

SCE’s induction and its relation with age groups

The results that presented in table demonstrate the effect of age groups (years)
on induction of SCEs in the brick kiln workers, showed that the average of SCE
frequency per cell increased with the increased ages of the workers of brick kiln.
The SCE frequency /cell at group age at (10-19 years) gave 7.6±0.83 and increased
to 10.47÷0.96 at age group between 40-49. The results showed significant between
age group (40-49y) and age group (10-19y ) and age group (20-29y) while there is
no significant differences between age group (30-39) and age group (40-49 years)
and there was no significant differences between age group (10-19) and age group
(20-29 ). This results indicate that the workers in the brick kiln at age groups
(30-39) and (40-49 years) were more sensitive in induction of SEC than the other
groups and this result indicate that DNA repair process more efficient in group age
(10-19years) and age group (20-29 years ) than the other groups. The results showed
significant in age groups of 30=39 and age group 0f 1-5years old while there was no
significant at age groups between 10-19 years old and age group of 20-29 years old.
Aging is associated with a declining ability to maintain tissue structure and function.
Replicative senescence contributes to aging in two ways, by limiting regenerative
processes that require cell turnover and through the disruptive effects that senescent
cells have on tissues(Coppe et al., 2010). These effects are likely to be amplified
in vivo where the declining pool of replicatively competent cells would be called
upon to divide more often to fulfill the needs of tissue regeneration. The result is
to deplete the remaining cellular reproductive reserves at an ever accelerating pace.
With these effects in mind, we argue that high T-SCE rates and their connection to
accelerated replicative senescence forge a plausible mechanistic link to premature
aging in Werner’s and Bloom’s syndromes(Hagelstrom, et al ., 2010) The results

Table 2: Relationship between age group (year) and SCE’s in lymphocytes of brick kiln workers
No. Age groups (years) SCE’s ± SD
Group 1 10-19 7.62 ± 0.83
Group 2 20-29 8.28 ± 0.79
Group 3 30-39 10.25 ± 1.04
Group 4 40-49 10.47 ± 0.96
Control Different ages 4.38±0,15
LSD value 2.168 *
P-value 0.0371
* (P<0.05)

presented in table -3 demonstrate the effect of employee years on induction of SCE

62
in the workers of brick kiln. The average of SCE frequency was in employee years
between 16-20 years that gave the highest average of SCE frequency /cell it gave
10.74 SCE’s per cell, while the employee years between 1-5years gave the lowest
value of SCE , it gave 7.67. The results showed significant on level of probability
0.05 between employee years 0f 21-25 years and 1-5 years and 6-10 years but
the results showed no significant in groups of 11-15 years of employee and 16-20
years by using LSD value. These results indicated that induction of SCEs increase
with the increase of employee years because of increase of pollutants of brick kiln
which exposed by occupational workers (Sears,2007). Also the results discuss the

Table 3: Relationship between service years and SCE’s in lymphocytes of brick kiln workers
No. Service years SCE’s ± SD
Group 1 1-5 7.69±0.061
Group 2 6-10 9.42±0.77
Group 3 11-15 10.46±0.62
Group 4 16-20 11.74±0.71
Group 5 21-25 10.55±0.60
Control 0 4.38 ±0.15
LSD value 1.773 *
P-value 0.0371
* (P<0.05)

relationship between the effect of smoking and employee years and its relation
with SCE frequency. The results demonstrate significantly by using LSDvalue at
probability 0.01 and 0.05, that the smokers have higher value of SCE with employee
rears at 16-20 years which gave 9.15 SEC /cell and the value of SEC was decreased
with group of employee rears between 1-5 years to give the lowest value ( 5.5
SEC /cell). Smoking and age service years are well known determinants of SCE
frequency, and they also significantly influenced the SCE frequency in our study.
After adjustment for these confounding variables, a significantly higher mean SCE
frequency was observed for the nonsmoking, this results agreed with Lander and
Rφnne2015 . The study of SCE’s is considered one of the genetic changes that

Table 4: Relationship between smoking and non-smoking with SCE’s in blood lymphocytes of brick
kiln workers
SCE’s ± SD
Service years LSD value P-value
Smoking Non-smoking
1-5 10.10ś0.54 5.50±0.17 2.319 ** 0.0138
6-10 11.82ś0.66 7.90±0.31 1.973 * 0.0352
11-15 12.38ś0.59 7.40 ±0.28 2.406 ** 0.0155
16-20 14.34ś0.81 9.15±0.47 2.317 ** 0.0038
21-25 12.80ś0.60 9.42±0.36 1.062 * 0.0471 *
Control 5.38±0.16 4.38±0.15
* (P<0.05), ** (P<0.01)

happen in the cells and it help to determine the mutagens and the pollutants in the
contaminated environment like the brick kiln in vivo and in vitro directly or after
metabolically activated especially by certain concentrations of pollutant that human
being exposed to it. The frequency of SEC in different cells give an idea of DNA

63
 damage. The SEC’s is affected by different pollutants that effect in cell cycles and
also changed by the time of replication. The SCE’s frequency show correlation
with the mutation of single gene. The increase in SCE’s frequencies is a test for
genetic analysis that present of two double helices of DNA, the two chromatids was
opened and then rejoined again in new loci, this rejoin does not happen in correct
manner, so it lead to fall of repair system in DNA and this will lead to increase of
the SCE’s frequencies which manifested by incorporation 0f 5-BuDR and examines
by fluorescent microscope by using DAPI stain (Laqqan, et al., 2012).

Conclusion

The results showed that the emitted gases and fumes gave significant effect on the
human health that are working and exposed to the pollutants of brick kiln. The data
demonstrate that the effect of the fumes and gases on workers by increasing the
frequencies of induction of SCEs which gave high value with the working smokers
and serving years. SCEs have been applied as a predictive assay of DNA damage
risk and as a biomarker of mutagenicity. Spontaneous or baseline SCE frequencies
in peripheral blood lymphocytes provide an additional diagnostic tool for identifying
accumulated chromosomal damage occurring during the lifespan of the cells.

References

Abid ,K.S and Salman, J.H.(2005). Study the reality of brick factories in Nahrawan
technical planning and fellow up department of industrial activity / Ministry
of Environment.
Clancy,S.and. Shaw, K...M.(2008) . DNA Deletion and Duplication and the Asso-
ciated Genetic Disorders.Nature Education 1(1):23-28. Coppe JP, Desprez PY,
Krtolica A, Campisi J. The senescence-associated secretory phenotype: the
dark side of tumor suppression. Annu Rev Pathol 2010; 5:99-118.
Jafar, Z.M.T., Shubber,E.K ., and Amash,H.M. (2001) . Cytogenetics analysis of
Chinese hamster lung fibroblast spontaneously resistant to methotrexate . The
Nucleus , 44:28-35.
Laqqan,M.M., Ashour,M.J., Silmi,A.S., and Sharif,A. (2011). Chromosomal Aber-
rations in males occupationally exposed to chemical pollutants in the Gaza
Strip Palestine . The Islamic university journal (Series of Natural studies and
engineering .19(1):23-31.
Hagelstrom RT, Blagoev KB, Niedernhofer LJ, Goodwin EH, Bailey SM. Hyper
telomere recombination accelerates replicative senescence and may promote
premature aging. Proc Natl Acad Sci U S A 2010; 107:15768-15773.
Lander F, Rφnne M 1995 . Frequency of sister chromatid exchange and hematolog-
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64
Rabin , D 2015. Work related Injuries and Musculoskeletal Disorders among Child
Workers in the Brick Kilns of Khejuri of Purba Mediipur in West BengalInter-
national Journal of Advanced Research (2015), Volume 3, Issue 3, 1065-1076.
Rudd MK, Friedman C, Parghi SS, Linardopoulou EV, Hsu L, Trask BJ (2007)
Elevated Rates of Sister Chromatid Exchange at Chromosome Ends. PLoS
Genet 3(2): e32. doi:10.1371/journal.pgen.0030032.
Saber ,A.T., Jacobsen, N.R., Jackson , P., Poulsen ,S.S., Kujovska, Z.O., Halap-
panavar,S., Yauk,C.L., Wallin , Hand Vogel , U. (2014). Particle- induced pul-
monary acute phase response may be the causal link between particle inhalation
and cardiovascular disease . Wiley Inter Rev. Nanomed.Nanobiotechnol. 6(6):
517-531.
Sears, M.R.(2007). Lung function decline in asthma . Eur.Respir.J. 30:411-413.
Wojeik,E., Andraszek,K., Ciszewska,M., and Smalec, E. (2012). Sister chromatid
exchange as an index of chromosome instability in chondrodystrophic chickens
.Gallus domesticus. 92:84-89.
Wilson III, D.M. and Thompson, L.H. (2007) Molecular Mechanisms of Sister
Chro-Matid Exchange. Mutation Research, 616, 11-23.
http://dx.doi.org/10.1016/j.mrfmmm.2006.11.017

65
66
 Nutritional Status in Relation to Intestinal Infection Among Primary
Schoolchildren in Erbil Province Kurdistan-Iraq
Abdullah A. Hama1,2 and Zohair I.F. Rahemo3
1
Department of Medical Laboratory , Technical College of Health, Sulaimani
Polytechnic University, Kurdistan-Iraq
2
Department of Medical Laboratory Science , Komar University for Science and
Technology, Kurdistan-Iraq
3
Department of Biology, College of Science, Salahaddin University, Kurdistan-Iraq
kany [email protected]
Abstract
In the present study the status of malnutrition (According to WFA) was 15.36%
and 8.56% among infected children with intestinal parasites and non-infected chil-
dren respectively. The malnutrition was more common among children who in-
fected with more than one parasite, the rate of malnutrition (WFA) from positive
cases were 44%, 62%, and 80% for single, double and triple parasitic infections
respectively. Intestinal parasitic infections were more prevalent in underweight chil-
dren (low WFA), whereas wasted children (low HFW) have slight associated with
parasitic infections, but there is no relation between stunted children (low HFA) and
intestinal parasitic infection. Also noted that H. nana, B. hominis and G. lamblia
were more common prevalence among children suffer from malnutrition.
Keywords: malnutrition, intestinal parasite, school age children

Introduction

Parasitic infection especially intestinal protozoan remains an important public health

concern worldwide (Olivares et al., 2002). Human intestinal parasitic is commonly
spread through the ingestion of food or water contaminated with human or animal
faeces as a result of poor sanitation (Bhandari et al. 2011; Katz and Taylor, 2001).
Intestinal parasitic infections are more common among school-age children and tend
to cause high intensity in this age group which leading nutritional deficiency and
impaired physical development is likely to have negative consequences for cognitive
function and learnability (Ulukanligil and Seyrek, 2003). The correlations between
malnutrition, intestinal infection, and child development are complex and studies of
these interrelationships will allow health agencies to maximise screening and inter-
vention strategies for developing countries (Rodriguez et al. 2011). Erbil province
consider as developing region in Kurdistan of Iraq. Although several studies indi-
cated that intestinal parasite infection was endemic among the community, there is
no adequate data about the intestinal parasite infection and their relations to nutri-
tional status among schoolchildren. In Iraq relatively there were few studies done
on intestinal parasites and to the best our knowledge, no studies have been done on
the intestinal parasites and their effect on nutritional status in the Kurdistan region.

67
 Materials and methods

Sampling

This study was achieved on 1028 children included 512 boys and 516 girls ranging
6-12 years old in both urban and rural areas randomly. All samples were taken from
nine primary schools in different areas. Stool samples were collected from all cases
in clean screw-cap containers with/or without preserved solution, for the cellophane
tape test mother was asked to perform cellulose tape test from her child/children at
morning before defecation (Okyay et al., 2004). The stool specimens were exam-
ined microscopically for the detection of intestinal parasites using both direct and
concentration methods

Anthropometric measurements

Weight measured to the nearest 50 g with a mechanical scale. Height measured

to the nearest 0.1 cm with a portable Stadiometer. Nutritional status evaluated
by calculating, weight-for-age (WFA), weight-for- height (WFH) and height-for-
age (HFA) according to National Center for Health Statistics (NCHS, 2000) chart.
Statistical analysis was carried out using software (SPSS version 17, chi-square test).

Results

Malnutrition determined as weight-for-age (WFA), weight-for-height (WFH), and

height-for-age (HFA) according to National Center for Health Statistics (NCHS,
2000). The general prevalence of malnutrition among schoolchildren was 23%
(WFA), 20.4% (HFA) and 6.7% (WFH). Intestinal parasite infection was more
prevalent in underweight children (low WFA), whereas in order wasted children
(low HFW) was slightly associated with parasitic infections, but there is no relation
between stunted (low HFA) children and parasitic infection.
G. lamblia, B. hominis and H. Nana were most prevalent in malnourished chil-
dren, and the rate of malnutrition was slightly higher in males (12.35%) than in
females (11.57%) in type WFA.
Rates of malnutrition among infected children with intestinal parasites were
15.36%, 8.36% and 3.5% for WFA, HFA and WFH respectively, versus non-infected
children, 8.56%, 12% and 3.3% for WFA, HFA and WFH respectively. In present
study higher rates of malnutrition (WFA) were found in the infected group (Table ).
As illustrate in Table the rates of malnutrition were more prevalent in children
who infected with more than one parasite (multiple infections). The percentage of
malnutrition among infected children for (WFA) were 44%, 62%, 80% of single,
double and triple infection respectively, and for (HFA) were 24.6%, 32.5%, 40% of
also respectively, for (WFH) were 9.5%, 18.6%, 20% of single ,double and triple in-
fection respectively, which clearly indicated the relation between multiple infection
and malnutrition.Important Remarks.

68
 Table 1: Malnutrition rates among infected and non-infected children.
Infected children
Non-infected Total rate of malnutrition
Malnutrition with intestinal
children From (1028)
parasites
No % No % No %
Male 83 8.07 44 4.28 127 12.35
WFA Female 75 7.29 44 4.28 119 11.57
Total 158 15.36 88 8.56 246 23.9
Male 39 3.79 48 4.66 87 8.46
HFA Female 47 4.57 76 7.39 123 11.96
Total 83 8.36 124 12.06 210 20.4
Male 19 1.84 20 1.94 39 3.79
WFH Female 17 1.65 14 1.36 31 3.01
Total 36 3.5 34 3.3 70 6.8
X 2 = 24.76, 2
P − value = 0.0001

Table 2: Malnutrition among infected children with single and multiple infections.
Infected children with No. of WFA HFA WFH
Sex
intestinal parasites positive cases No % No % No %
Male 64 22.5 42 14.8 13 4.6
Single infection 284 Female 63 22.1 28 9.8 14 4.9
Total 127 44.7 70 24.6 276 9.5
Male 16 37.2 8 18.6 5 11.6
Double infection 43 Female 11 25.5 6 13.9 3 6.9
Total 27 62.7 14 32.5 8 18.6
Male 3 60 2 40 1 20
Triple infection 5 Female 1 20 0 0 0 0
Total 4 80 2 40 1 20
Male 83 25 52 15.6 19 5.7
Total infection 332 Female 75 22.5 34 10.2 17 5.1
Total 158 47.5 86 25.9 36 10.8
WFA= Weight for Age; HFA= Height for Age; HFW= Height for Weight

Discussion

The malnutrition (WFA) has been related strongly with intestinal parasite infection
because they produce many intestine health problems, such as intestinal dysmotility,
malabsorption, constipation, epithelial invasion, covering of an intestine surface
area, blood loss, diarrhoea (steatorrhea), all these factors may contribute to weight
loss and malnutrition.
Intestinal parasitic infection also causes nutritional deficiency and impaired phys-
ical growth (Ulukanligil and Seyrek, 2003) also affect childhood development and
morbidity in many countries, on the other hand, malnutrition may increase sus-
ceptibility to parasitic infection (Quihui-Cota et al., 2004). In the present study,
we observed a clear association between the intestine infection with parasites and
malnutrition, this result supported by several studies which indicated that malnu-

69
 trition is one of the most common sequences of infection with intestinal parasites.
The malnutrition has been found to be associated with school achievement in many
countries (Hutchinson et al., 1997). Oberhelman et al. (1998) in Nigeria agreed
with the current study they concluded that parasites were more prevalent in children
with low (HFA) and (WFA), also Tsuyuoka et al. (1999) in Brazil; Ulukanligil and
Serkey (2004) in Turkey have reported that the children with parasites had lower nu-
tritional indices, and they found significant association between intestinal helminths
and malnutrition. Also, this finding agreed with (Hama and Rahimo 2014) they
state that intestinal parasite has a direct effect on the haemoglobin concentration
and physical growth.
Some studies didn’t agree with the current study De-Morais et al. (1996) in
Brazil, Awasthi and Pande (1997) in Lucknow, India, Al-Agha and Teodorscu (2000)
in Palestine, Muniz et al.(2002) in Brazil, and Rai et al. (2004) in Nepal, these may
be due to high standard of socio-economic level and the best quality of meals of
the studied area, or intestinal parasites intensity and type of infection have notable
role reveals the symptoms as well as signs of malnutrition.
It is noteworthy to report that some studies were conducted among acute infection
(not chronically) as such in patients, we cannot correlate the effect of parasites on
nutritional status clearly because the period of infection is not efficient to assess
the nutritional parameters change, but chronic and asymptomatic intestinal parasitic
infections cause nutrition deficiency.

References

Al-Agha, R. and Teodorescu, I. (2000). Intestinal parasites infestation and anemia

in primary schoolchildren in Gaza province, Palestine. Rum. Arch. Microbiol.
Immunol., 59 (1-2): 131-143.
Awasthi, S., and Pande, V. K. (1997). Prevalence of malnutrition and intestinal
parasites in slum preschool children in Lucknow. Ind. J. Pediater., 34 (7):
599-605.
Bhandari N1, Kausaph V and Neupane GP. (2011). Intestinal parasitic infection
among school-age children. J Nepal Health Res Counc. (1):30-2.
De-Morais, M.B.; Suzuki, H. U.; Coral, J. N.; Machado, N.L., and Neto, N.F.
(1996). Asymptomatic giardiasis does not affect iron absorption in children
with iron deficiency anaemia. Med. J., 15 (5): 421.
Hama, A. A. and Rahimo. Z.I (2014). Intestinal Parasitosis in Relation To
Haemoglobi Concentration Among school age children in Erbil province Kur-
distan, Iraq, International Science Journal, 1(1), 96Ű99.
Hutchinson, S.E.; Powell, C.A.; Walker, S.P.; Chang, S.M., and Grantham-Mc-
Gregor, S.M. (1997). Nutrition, anaemia, geohelminth infection and school
achievement in rural Jamaican primary schoolchildren. European J. Clin. Nut.,
51: 729-735.

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Katz, D.E. and Taylor, D.N. (2001). Parasitic infections of the gastrointestinal tract.
Gastroenterol. Clin., 30 (3): 1-31.
Muniz, P.T.; Ferreira, M.U.; Ferreira, C.S.; Conde, W.L., and Monteiro, C.A.
(2002). Intestinal parasitic infections in young children in Sao Paulo, Brazil:
prevalence, temporal trends and associations with physical growth. Ann. Trop.
Med. Parasitol., 96 (5): 503-512.
NCHS (National Center for Health Statistics)(2000).Growth charts.
http://www.cdc.gov/growthcharts.
Oberhelman, R.A.; Guerrero, E.S.; Fernandez, M.L.; Silio, M.; Mercado, D.;
Comiskey, N.; Ihenacho, G., and Mera, R. (1998). Correlation between in-
testinal parasitosis, physical growth and psychomotor development among in-
fant and children from rural Nicaragua. Amer. J. Trop. Med. Hyg. 58 (4):
470-475.
Okyay, P.; Ertug, S.; Gultekin, B.; Onen, O., and Beser, E. (2004). Intestinal
parasites prevalence and related factor in schoolchildren a western city sample-
Turkey. J. Bio. Med. Cen. Publ. Hlth., 22 (4): 64.
Olivares, J.L.; Fern?ndez, R.; Fleta, J.; Ruiz, M.Y. and Clavel, A. (2002). Vitamin
B12 and folic acid in children with intestinal parasitic infection. Amer. J., 21
(2): 109-113.
Quihui-Cota, L..; Valencia, M.E.; Crompton, D.W.; Phillips, S.; Hagan,P.; Cama-
cho,S.P, and Tejas,A.T. (2004). Prevalence and intensity of intestinal parasitic
infections in relation to nutritional status in Mexican schoolchildren. J. R. Soc.
Trop. Med. Hyg., 98 (11): 653-659.
Rai, S.K.; Hirai, K.; Abe, A.; Nakanish, M.; Rai, G.; Uga, S., and Shrestha,
H.G.(2004). Study on enteric parasitosis and nutritional status of school chil-
dren in remote hilly areas in Nepal. Nepal, Med. Coll. J., 6 (1): 1-6.
Rodr?guez, L., Cervantes, E., and Ortiz, R. (2011). Malnutrition and Gastrointesti-
nal and Respiratory Infections in Children: A Public Health Problem. Interna-
tional Journal of Environmental Research and Public Health, 8(4), 1174Ű1205.
doi.org/10.3390/ijerph8041174
Tsuyuoka, R.; Bailey, J.W.; Nery, G.A.M.; Gurgel, R.Q. and Cuevas, L.E. (1999).
Anemia and intestinal parasitic infections in primary school students in Aracaju,
Sergipe, Brazil. Cad. Saude, Publ., 15 (2): 413-421.
Ulukanligil, M. and Seyrek, A. (2003). Demographic and parasitic infection status
of schoolchildren and sanitary conditions of schools in Sanliurfa, Turkey. BMC
(Bio. Med. Central) Publ. Hlth., 118 (2): 151-158.
Ulukanligil, M. and Seyrek, A. (2004). Anthropometric status, anaemia and in-
testinal helminthic infections in shantytown and apartment schoolchildren in
the Sanliurfa province of Turkey. Eur. J. Clin. Nutr., 58 (7): 1056-1061.

71
72
Chemistry Session
Evaluation of Hassira Crude Oil in Sarqalah, Kifri, Sulaimani, Kurdistan
Region
Hersh J. Noori and Nabil A. Fakhre
Department of Chemistry, Garmian University, Kallar, Sulaimani, Kurdistan-Iraq
[email protected]
Abstract
Evaluation of Hassira (Ha) crude oil in Sarqalah, Kifri, Sulaimani, and Kurdistan
has been studied. The assay provides general data on the oil and was based on some
tests such as distillation range, water content, S.G., API, R.V.P., flash point, pour
point, ash content, sulphur content and elemental analysis that enable desirable or
undesirable features to be noted. As a result, according to results obtained for Has-
sira field in the Hassira Village, Sarqalah, Kifri, Sulaimani in Kurdistan Region and
evaluated with world specification, it has a high level of API is more than 40 thus
indicate it is light crude oil means the aliphatic more than aromatic components.
In addition we compared with other Kurdistan crude oil, (Ha) after Taq-Taq 1 is
the best type of crude oil in Kurdistan such as (Taq-Taq2, Tawki1, By-Hassan and
Kirkuk).
Keywords: Crude oil, Aliphatic and aromatic compounds, elemental analysis

Abbreviations:
SG, Specific gravity; API, American petroleum institute; RVP, Ried vapour pressure;
EIEB, extraction induced emulsion breaking; ICP-OES, inductively coupled plasma
optical emission spectrometry.

Introduction

Most crude oils form over millions of years from the remains of tiny aquatic plants
and animals that are exposed to the combined effects of time and temperature. It is
a complex mixture containing both organic and inorganic chemical species, being
trace metals one group the inorganic components present (Hardaway, Sneddon &
Beck 2004, p.2883). Crude oils mainly consist of hydrocarbons mixed with variable
amounts of sulphur, nitrogen, oxygen compounds. Three general classes of hydro-
carbons arc present in crude oil paraffins, naphthenes, and aromatics (Chaudhuri
2016, p.1). The petroleum industry generally classifies crude oil by the geographic
sources it is produced in, its API gravity, and its sulfur content. The Carpatica
classification divided crude oils into seven classes based on the percentage of paraf-
finic, naphthenic, and aromatic carbon compounds (Creanga 1961, p.90). The metals
present in the crude oils are mostly Ni (II) and V (II) porphyrins and non-porphyrins.
Other metal ions reported from crude oils; include Copper, Lead, Iron, Magnesium,
Sodium, Molybdenum, Zinc, Cadmium, Titanium, Manganese, Chromium, Cobalt,
Antimony, Uranium, Aluminum, Tin, Barium, Gallium Silver and Arsenic. Metel-
loporphyrins are among the first compounds identified belonging to the biological
origin (Treibs 1936, p. 683).

73
 Information on trace element concentrations in crude oil is getting increasingly
important for the geochemical characterization of source rocks and basins and for
corrective actions during oil production and refining (e.g. prevention of scale forma-
tion and catalyst poisoning, corrosion and pollution control)( Sainbayar, Monkhoo-
bor & Avid 2011, p.22). The heavy metal contents of crude oils are often correlated
with percentages of light, middle, heavy and residual fractions of crude oils. For
heavy crude oil upgrading, the high metals levels necessitate the processing to be
costly, either through the high costs of hydrogen addition or high yield loss due to
carbon rejection (Barbooti 2015, p. 326). This process in turn requires the availabil-
ity of reliable method for metal determination. The techniques used most frequently
for the determination of trace elements in crude oil are inductively coupled plasma
optical emission spectrometry (ICP OES)( De Oliveira, Ribeiro & Carneiro 2015,
p.182), (Tugsuu, Yoshikazu & Enkhsaruul 2012, p. 402) & (Karim, Hamasalih
& Tofiq 2015, p.75) atomic absorption spectrophotometry, F-AAS (Barbooti 2015,
p. 325) & (Pessoa, Hauser-Davis & de Campos 2012, p. 1569), electrothermal
atomic absorption spectroscopy (ET-AAS) (Duyck, Miekeley & da Silveira 2007,
p.940) and high performance liquid chromatography hplc (Tadayon, Massoumi &
Eslami 1999, p.372). The aim of this study to Evaluation of Hassira (Ha) crude
oil in Sarqalah, Kifri, Sulaimani, Kurdistan Region then comparisons with some
other Kurdistan fields of such as Taq Taq 1 (Tq1), Taq Taq 2 (Tq2), Tawki 1(Tw1),
Hamreen (Hr), Bay Hassan (Bh) and Karkuk (Kru).

Materials and Methods

Material and Reagents

Hassira crude oil provided form western zagros production field in Sarqalah, Kifri,
Sulaimani, Kurdistan Region.

Procedures

Sample Preparation for trace elements determination

The elements were extracted from crude oil sample through the extraction induced
by emulsion breaking (EIEB). The procedure is based on the formation of stable
emulsions between the samples and an acid Triton X-100. To improving the ex-
traction of the metals, in the initial, 80% of the samples (approximately 5.2 g) and
20% of the 3% the acidic Triton X-100 solution were mixed and stable water-in-oil
emulsions were obtaining by vigorous mixing of the samples. The emulsions were
agitated for 30 min on a mixer (auto vortex). To promote the emulsion breaking
the mixed solution was putting in the centrifuge at 3000 rpm for 30 minute. The
emulsion breaking resulted in two well-separated phases: an upper phase contain-
ing the organic part of the solution and a lower phase containing the aqueous acid
solution with the extracted elements. Then, the upper phase was removed and the
lower aqueous phase was collected and transferred the calibrated flasks then diluted

74
to the mark with deionized water. Finally, the solution was used to measure the
elements of interest.

Evaluation tests for crude oil

Specific gravity and API gravity

According to ASTM D-1298 (Hydrometer test-The sample [crude oil] was brought
to the prescribed temperature transferred to a cylinder at approximately the same
temperature. The appropriate hydrometer was lowered into the sample and allowed
to settle. After temperature controlled, the hydrometer scale was read and the
temperature of the sample was noted.
Flash point
According to ASTM D-92 (Cleveland open cup-The test cup was filled to a specified
level with the sample [non-volatile petroleum fractions]. The temperature of the
sample increased rapidly at first and then at a slow constant rate as the flash point is
approached. At specified interval a small test flam was passed across the cup. The
lowest temperature at which application of the test flame caused the vapor above
the surface of the liquid to ignite was taken as the flash point.
Ried vapour pressure (RVP)
According to ASTM D 323 a chamber of the vapour pressure apparatus was filled
with the chilled sample and connected to the vapour chamber that has been heated
to 37.8◦ C (100◦ F) in a bath. The assembled apparatus was immersed in a bath at
37.8◦ C (100◦ F) until a constant pressure was observed.
Water and sediment content
According to ASTM D2709 a 10-mL sample of the undiluted crude oil was cen-
trifuged at a relative centrifugal force of 800 (rpm) for 10 min. at 21 to 32◦ C (70
to 90◦ F) in a centrifuge tube readable to 5 mL and measurable to 1 ml. After
centrifugation, the volume of water and sediment which was settled into the tip of
the centrifuge tube was read to the nearest 5 mL and reported as the volumetric
percent water and sediment by centrifuge.
Salt content
The test method was based on the conductivity of a solution of samples in a polar
solvent when subjected to an alternating electrical stress. According to ASTM D
3230 a sample was dissolved in a mixed solvent and placed in a test cell consisting
of a beaker and two parallel stainless steel plates. An alternating voltage was passed
through the plates and the resulting current flow was shown on a milliammeter. The
salt content is obtained by reference to a calibration curve of current versus salt
content of known mixtures.

Determination of Total Sulfur Content

By XRF-Technique A 5g sample of crude oil were introduced into the sample cell of
XRF instrument which was calibrated before sample measuring by standards from
(0.0 Ű 1.0% and 1.0 Ű 6.0%) sulfur .The sample was put in a sample holder of the
instrument and the total sulfur content in the sample was measured and recorded.

75
 Result and Discussion

Evaluation of crude Oils Produced According to the ASTM Standards

Crude oil exhibits wide variations in composition and properties. The assay provides
general data on the oil and was based on some tests such as water content, S.G., API,
R.V.P., flash point, pour point, ash content, sulphur content and elemental analysis
that enable desirable or undesirable features to be noted (Miguel, Da Cunha &
Pacheco 2010, p.1506). Table (1) shows these properties and specification of Ha
crude oil, TQ1and TQ2 (Karim, Hamasalih & Tofiq 2015, p.75) and Tu, Bh, Hr and
Kru (Safea 2007) in comparison to world requirement for marketing specification
of crude oil.
Density or specific gravity is used for conversion between mass and volume
measurements, and it is used also for determination of API gravity and with other
test results to predict oil quality. With increasing API quality of crude oil will
increase thus the prices of it will raise (Kaes 2000, p. 12). As it is showed in table
(3.1) the TQ2 crude oil has the lowest degree of API which is about 23.74, then
each of Tw, Bh, Hr and Kru has API of 29.8, 32.19, 32.19 and 36.05 respectively.
While API of Ha crude oil is 40.2 and TQ1 is 47.52. As a result, TQ1 has the
maximum API level then Ha. This means that the TQ1 and Ha have the highest
percent of paraffins and lowest aromatics and contains a lowest ratio of sulphur,
nitrogen, carbon residue and asphalting as illustrated in fig. (1).

Figure 1: Illustration of the general relationship of petroleum (Kaes 2000, p. 13)

Sediments are tiny solid particles; also they have effect on the quality of crude
oil. According to world specification, the normal range of sediment in crude oil is
(0.01-0.40%) (ASTM D2709 2004). It is noticed that the sediment percent of Ha
is 0.02%, which is the lowest level, and for each of TQ1, TQ2, TW and Kru is
0.05, 0.08, 0.227, and 0.2 respectively. Ha crude oil sample has high value of salt
content (0.245 ppm) in compare with the other samples except Kru sample.
F.P. and RVP have vital role in determining the ratio of gas and light materials
in oil. Moreover, they can be applied approximately to evaluate the heaviness and

76
lightness of oil (ASTM D92 2004) & (ASTM D323 2004). Table (1) shows the
flash point of the samples, in which (Ha, TQ1, TQ2, Kru) nearly have the same
flash point except TW1 crude oil sampler which has high flash point (40 ◦ C). The
RVP of (Ha) sample was (31 Kpas) increase with increasing volatile materials such
as natural gases.
The normal range of sulphur content is between 0.1-5 percent. In terms of total
sulphur content (Ha) crude oil is considered one of the best oils around Kurdistan,
in which sulphur content of TQ1 0.61%, while Tw oil has maximum sulphur content
2.7% which is still within the normal range (ASTM D4294 2004).

Table 1: Physicochemical properties of some crude oil fields in Kurdistan.

Test description Ha TQ1 TQ2 Tw Bh Hr Kru

Sp. Gr. at 15.5◦ C 0.8241 0.7904 0.912 0.88 0.902 0.864 0.844
API Gravity 40.201 47.52 23.74 29.08 32.19 32.19 36.05
Water content (v/v%) 0 0 0 0.09 2.5 5.6 0.00
Sediments (%) 0.02 ¡ 0.05 0.08 0.227 — — 0.2
Salt content (ppm) 0.245 0.0024 0.02 0.01 0.18 0.004 5.0
Pour point (◦ C) — - 40 - 27 -30 -20 -25 -25
Ash content (wt.%) 0.023 0.047 0.058 0.015 0.26 0.034 0.009
Flash point ◦ C (C.O.C) Flam-mable Flam-mable Flam-mable 40 — — Flam-mable
Water and sediment (V%) 0.02 ¡ 0.05 0.08 0.317 — — 0.2
Total sulphur content (mass%) 0.9838 0.61 2.20 2.70 — — 1.99

Metal Analysis in Crude oil

The presence and concentration of trace elements vary with the geographical source
of a crude oil. Even minute amounts of iron, copper, and particularly nickel and
vanadium in the charging stocks for catalytic cracking affect the activity of the cat-
alyst and result in increased gas and coke formation and reduced yields of gasoline.
It was found that cadmium, zinc, iron and copper were not indigenous to the crude
oil but were the result of contamination with associated water and/or particles from
the producing wells (Tugsuu, Yoshikazu & Enkhsaruul 2012). According to the
table (2) the concentration of the trace elements in Hassira crude oil is very low,
except of Ca, Na and S which have the values (16184, 27856 and 31008) ţg/L using
ICP-OES.

Table 2: Determination the elements in (Ha) crude oil applying EIEB using ICP-OES in (ţg/L)
Elements Determined conc. Elements Determined conc. Elements Determined conc.
Ca 16184 Be 30 Ce 0.485
Li 516 Si ¡DL Fe 2611
Mn 155 Ti 430 Cu 138
Na 27855 Zr 263 V 5159
K 2447 Zn 15600 Cd ¡DL
Ni 6988 As 1580 Cr ¡DL
Pb ¡DL Mo 41

77
 Conclusion

It was concluded that Hassira crude oil field has a high level of API is more than 40
thus indicate it is light crude oil means the aliphatic more than aromatic components.
According to results obtained and evaluated with world specification, we compared
with other Kurdistan crude oil, (Ha) after Taq-Taq 1 is the best type of crude oil in
Kurdistan such as (Taq-Taq2, Tawki1, By-Hassan and Kirkuk).

References

ASTM D2709, 2004. Standard Test Method for Water and Sediment in Middle
Distillate Fuels by Centrifuge.
ASTM D323, 2004 An American National Standard Test Method for Vapor Pressure
of Petroleum Products (Reid Method).
ASTM D4294, 2004 Standard Test Method for Sulfur in Petroleum and Petroleum
Products by Energy-Dispersive X-ray Fluorescence Spectrometry.
ASTM D92, 2004. Standard test method for Flash point by Cleveland open cup.
Barbooti, M.M., 2015. Evaluation of Analytical Procedures in the Determination of
Trace Metals in Heavy Crude Oils by Flame Atomic Absorption Spectropho-
tometry. American Journal of Analytical Chemistry, 6(04), p.325. 325.þ
Chaudhuri, U.R., 2016. Fundamentals of petroleum and petrochemical engineering.
CRC Press.
Creanga, C., 1961. Chemical classification of crude oils. Carpatica classification.
Acad. Rep.
de Oliveira Souza, M., Ribeiro, M.A., Carneiro, M.T.W.D., Athayde, G.P.B., de
Castro, E.V.R., da Silva, F.L.F., Matos, W.O. and de Queiroz Ferreira, R.,
2015. Evaluation and determination of chloride in crude oil based on the
counterions Na, Ca, Mg, Sr and Fe, quantified via ICP-OES in the crude oil
aqueous extract. Fuel, 154, pp.181-187.
Duyck, C., Miekeley, N., da Silveira, C.L.P., Aucelio, R.Q., Campos, R.C., Grin-
berg, P. and Brandao, G.P., 2007. The determination of trace elements in crude
oil and its heavy fractions by atomic spectrometry. Spectrochimica Acta Part
B: Atomic Spectroscopy, 62(9), pp.939-951.
Hardaway, C., Sneddon, J. and Beck, J.N., 2004. Determination of metals in crude
oil by atomic spectroscopy. Analytical letters, 37(14), pp.2881-2899.
Kaes, G.L., 2000. Refinery Process Modelling, Kaes Enterprises. Inc. ABD.
Karim, A.R., Hamasalih, L.O., Tofiq, D.I. and Salih, D.M., 2015. A Comparative
Elemental Analysis of the Different Kurdistan Crude Oil Fields. Technology,
2(2), pp.74-78.

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Miguel, E.M., Da Cunha, A.L., Pacheco, W.F., Farias, P.A. and Aucélio, R.Q.,
2010. Bismuth film anodic stripping voltammetry for the determination of
lead in kerosene: a metrological study. Electroanalysis, 22(13), pp.1505-1510.
Pessoa, H.M., Hauser-Davis, R.A., de Campos, R.C., de Castro, E.V.R., Carneiro,
M.T.W.D. and Brandäo, G.P., 2012. Determination of Ca, Mg, Sr and Ba in
crude oil samples by atomic absorption spectrometry. Journal of Analytical
Atomic Spectrometry, 27(9), pp.1568-1573.
Safea, S. T., 2007. Evaluation of crude oil and its products of Tawke Well in
Zhakho-Kurdistan Region and gasoline octane number improvement by some
additives. Ph.D. Thesis, Salahaddin University.
Sainbayar, J., Monkhoobor, D. and Avid, B., 2011. Determination of trace elements
in the Tamsagbulag and Tagaan Els crude oils and their distillation fractions
using by ICP-OES.
Tadayon, F., Massoumi, A. and Eslami, M., 1999. Determination of vanadium,
nickel, and iron in crude oil by high-performance liquid chromatography. Jour-
nal of chromatographic science, 37(10), pp.371-374.
Treibs, A., 1936. Chlorophyll and minerals in organic Häminderivate. Angew
Chemie, 49, pp.682-686.
Tugsuu, T., Yoshikazu, S., Enkhsaruul, B. and Monkhoobor, D., 2012. A compar-
ative study on the hydrocracking for atmospheric residue of Mongolian Tam-
sagbulag crude oil and other crude oils. Advances in Chemical Engineering
and Science, 2(03), p.402.
Tugsuu, T., Yoshikazu, S., Enkhsaruul, B. and Monkhoobor, D., 2012. A compar-
ative study on the hydrocracking for atmospheric residue of Mongolian Tam-
sagbulag crude oil and other crude oils. Advances in Chemical Engineering
and Science, 2(03), p.402.

79
80
Batch and Flow Injection Spectrophotometric Determination of Amoxicillin in
Pharmaceutical Formulations
Kamal M. Mahmoud* and Bery M. Rahman
Salahaddin University-Hawler (SUH); College of Science; Department of
Chemistry; Iraqi Kurdistan Region - Erbil
[email protected]
Abstract
This paper involves a batch and flow injection spectrophotometric determination
of amoxicillin (AMX) in pharmaceutical formulations. The methods were based on
the coupling reaction of AMX with diazotized sulphanilic acid in basic medium to
form a yellow color azo dye, which gives the maximum absorption at 439.5 nm.
In the batch method Beer’s law was obeyed in the concentration range of (0.1 -
11.0 µgml−1 ) and detection limit of 0.07 µgml−1 with a correlation coefficient (r)
of 0.9995 and a molar absorptivity of 2.9442×104 L mol−1 cm−1 , while the flow
injection analysis (FIA) system was linear under the optimum conditions in the con-
centration range of (1.0 - 35.0 µgml−1 ) with a detection limit of 0.7 µgml−1 , a
correlation coefficient (r) of 0.9970 and a sampling frequency 43S h−1 . The pre-
cision and accuracy of both methods were checked by calculating relative standard
deviation (RSD) and relative error (E%) for two different levels of concentration.
Possible interferences that related to the determination of AMX in pharmaceutical
formulations were studied and the results showed that the method was successfully
applied for determination of AMX in pharmaceutical formulations.
Keywords: Spectrophotometry, Flow Injection Analysis, Diazotization, Sulphanilic
acid.

Introduction

Analytical absorption spectroscopy in the Ultraviolet (U.V.) and visible region of the
electromagnetic spectrum has been widely used in pharmaceutical and biomedical
analysis for quantitative purpose and for the characterization of drugs, impurities,
metabolites and related substances [1].
Amoxicillin (AMX), or hydroxyl ampicillin was present as amoxicillin trihydrate
and amoxicillin sodium; chemically amoxicillin sodium was sodium (2S,5R,6R) -6-
[[(2R)-2- amino -2- (4- hydroxyphenyl) acetyl] amino] -3,3- dimethyl -7- oxo -4-
thia -1- azabicyclo [3.2.0] heptane-2-carboxylate. It was a semi-synthetic product
derived from a fermentation product, white, very hygroscopic powder, very soluble
in water, sparingly soluble in anhydrous ethanol, very slightly soluble in acetone
[2].
Various direct and indirect batch spectrophotometric methods [3-5] and flow in-
jection analysis (FIA) spectrophotometric have been reported for the determination
of AMX in pharmaceutical formulations [6-8]. A FI-CL method was described
for the determination of AMX [9-10]. Other methods were developed for the de-
termination of AMX such as spectrofluorimetric [11], potentiometric [12], atomic
absorption spectroscopy [13], HPLC [14] and HPLC-MS [15-16].
In this work selective, sensitive and accurate method was used for estimation of

81
 AMX by both batch and FIA. The method was based on the coupling reaction of
AMX with diazotized sulphanilic acid which formed from the reaction of sulphanilic
acid with nitrous acid in alkaline medium producing a water soluble yellow color
dye which was stable and has a maximum absorption at λmax 439.5 nm.

Experimental

Apparatus

The spectral and absorbance measurements were carried out on a (CECIL CE

3021 England) UV-Vis spectrophotometer, equipped with quartz cell of 1.0 cm
path length.
The schematic diagram of the flow injection spectrophotometric system used
in this work consists of a peristaltic pump (Carter 12/6 cassette pump with five
channels, variable speed) to deliver flow streams. The tygon pump tubes with (1.05
mm i.d.) were used to transport the carrier solutions. The rotary valve (Rheodyne
U.S.A.) with variable sample volumes was used. A Y-shaped Perspex piece was
used for mixing the two streams of reagents.
The absorbance measurements of the FIA were carried out using spectropho-
tometer (JENWAY 6405 UV/Vis Spectrophotometer) with flow cell Starna-micro
(100 µL) with 1.0 cm optical path length. The absorbance was recorded by means
of x-t recorder (type PM825A Philips, one line recorder) with various amplification
factors and different chart speeds.

Reagents and solutions

All chemicals and reagents were of analytical grade. Distilled water was used for
the preparation of all solutions.

Stock solution of amoxicillin

A stock solution (1000 µgml−1 ) of pure AMX was prepared by dissolving 1.0 g of
AMX sodium (provided by S.D.I.) in distilled water and diluted to 1.0 L. Working
standard solutions were prepared daily by appropriate dilution of the stock solution
with distilled water.

Sodium nitrite solution (0.2 M)

A fresh stock solution of 0.2 M was prepared daily by dissolving 1.38 g of sodium
nitrite (Fluka) in 100 ml of distilled water. Working standard solutions prepared by
appropriate dilutions of distilled water.

82
Sulphanilic acid (0.01 M)

0.01 M sulphanilic acid was prepared by dissolving of 0.866 g of sulphanilic acid

(BDH) in distilled water; the volume was completed to 500 ml. Working solutions
were prepared by suitable dilution with distilled water.

Sodium carbonate solution (0.5 M)

0.5 M of sodium carbonate solution was prepared by dissolving 26.5 g of Na2CO3

(Riedel-dehaën) in distilled water, then completed to 500 ml. Working solutions
were prepared by appropriate dilution.

Solutions (1.0 M) of sodium hydroxide, potassium hydroxide, ammonium hy-

droxide, hydrochloric acid, nitric acid, sulfuric acid, and acetic acid were prepared
each alone.

Sample preparation

Table (1) illustrates the trade products that contain AMX as pharmaceutical prepa-
rations, which were subjected to analysis by the following recommended method.

Ten tablets or the contents of ten capsules were weighed and grounded into a fine
powder (or mixed content of 5 vials). Powder equivalent to 50 mg of pure AMX
was dissolved in distilled water; the solution was stirred for 15 min, transferred to a
100 ml volumetric flask and completed with distilled water to obtain 500 µgml−1 .
The flask with its contents was shacked well and filtered (S&S Rundfilter). Working
solutions were prepared by diluting appropriate amounts of the stock solution with
distilled water.

Table 1: The trade name and composition of AMX contained in pharmaceutical formulations
Total amounts (mg per
Formulation Trade name Company Composition capsule or tablet or injec-
tion)
Amoxicillin
Amoxicillin S.D.I. (Iraq) 250
trihydrate
Capsule
Gobalpharma Amoxicillin
Glomox 500
(U.A.E.) trihydrate
HIKMA Amoxicillin
Penamox 500
(Jordan) trihydrate
HEUMANN Amoxicillin
Tablet Amoxicillin 500
(Germany) trihydrate
Acamoxil ACAI (Iraq) Amoxicillin 500
Injection sodium
Largopen Bilim (Turkey) Amoxicillin 500
sodium

83
 Solutions of interfering species

A stock solution of each interfering species (glucose, galactose, fructose, sucrose,

lactose, starch and sodium chloride) (BDH) was prepared by dissolving 1.0 g in
100 ml distilled water to obtained 10000 µgml−1 . Other solutions were prepared
by serial dilutions of the stock solutions.

The suggested mechanism

The possible mechanism for the reaction of diazotized sulphanilic acid with AMX
is explained in Fig. (1), the first step is formation of the diazonium ion and the
second step was coupling of diazonium ion with AMX in alkaline medium (sodium
carbonate) to form a yellow color complex which is soluble in water.

Figure 1: Possible mechanism of the reaction

Results and Discussion

Spectrophotometric determination of AMX using sulphanilic acid as diazotizing

agent

Preliminary work: To a 25 ml volumetric flask containing 2.0 ml of (0.1 M) hy-

drochloric acid solution, 2.0 ml of 0.01 M of NaNO2 solution, and 2.0 ml of 0.001 M
sulphanilic acid solution, 3.0 ml of aqueous sample solution containing 50 µgml−1
of AMX (6.0 µgml−1 AMX in 25 ml), and 2.0 ml of 0.25 M of Na2 CO3 solu-
tion were added, then the mixture diluted to the mark with distilled water. The
reagent blank was prepared in the same manner without AMX. These solutions
were prepared at room temperature, after 3 min the absorbance was measured at
439.5 nm.

84
Optimizations

To obtain the optimum conditions for the method, the following parameters were
investigated.

Absorption spectra

The conditions of preliminary work were used to obtain the absorption spectra
shown in the Fig. (2). Absorption spectra in visible region of the formed yellow
dye solution were scanned at different wavelengths between 300 - 550 nm. The
yellow color solution showed a maximum absorption at 439.5 nm (spectrum a),
while spectrum b for the blank solution showed zero absorbance at this wavelength,
therefore, this wavelength was selected for the spectrophotometric determination of
AMX.
It was known that the diazonium ion was prepared in acidic medium, therefore
the effect of acidity on the diazotization was studied; different types of acid solutions
(2 ml of 0.1 M) (hydrochloric acid, acetic acid, nitric acid, sulfuric acid) has been
tested. The results shown that there were no significant differences between the
types of acids on the absorbance of the formed azo dye, so as 0.1 M hydrochloric
acid was selected in the subsequent experiments [17-18].

Figure 2: Absorption spectra of (a) formed yellow azo dye against blank and (b) blank against
distilled water.

Effect of hydrochloric acid

Effects of different volumes (0.2 - 5.0 ml) of hydrochloric acid solution (0.1 M)
on the color intensity of the azo dye formed were investigated. Results obtained
maximum diazotization with 1.5 - 3.0 ml of 0.1 M hydrochloric acid, after that the
intensity of the color decreases with increasing the volume of hydrochloric acid,
because higher than this volume of hydrochloric acid concentration weakens the
alkaline medium required for obtaining the color of azo dye and at lower than this
volume the production of nitrous acid was not enough to form intense color of the
azo dye. Therefore, 1.5 ml HCl was used for subsequent work.

85
 Effect of sodium nitrite

Sodium nitrite was used for the formation of nitrous acid. Thus, the Effects of
different volumes (0.2 - 5.0 ml) of sodium nitrite (0.01 M) were tested for the
determination of AMX. Results indicated that 2.0 ml of NaNO2 solution gives
maximum absorbance of the colored product. Generally, sodium nitrite should not
be used in excess since an excess of nitrous acid exerts a very unfavorable influence
on the stability of azo solutions [19].

Effect of sulphanilic acid

The Effect of sulphanilic acid on the formation of the azo dye was examined in the
range of (0.5 - 10.0 ml); results indicates that (3.0 ml) of 0.001 M sulphanilic acid
exhibits maximum response. Higher than this volume an increase in the absorbance
were not significant, because this volume was enough to react with the nitrous acid
to form the diazonium ion.

Effect of the type of base

The preliminary experiments were shown that the colored dye can be developed
only in alkaline medium; therefore, the Effects of different alkaline solutions (0.25
M) were studied such as (sodium carbonate, sodium hydroxide, potassium hydroxide
and ammonium hydroxide). During these investigations, it was found that sodium
carbonate was the best suitable alkaline medium to obtain maximum absorbance;
therefore sodium carbonate was used in all subsequent experiments.

Effect of sodium carbonate

Effects of different volumes of sodium carbonate (0.25 M) were investigated in the

range of (0.5 - 5.0 ml). It was found that (2.0 ml) of the base gave the maximum
absorption intensity. Lower than this volume caused incomplete coupling reaction
between AMX and diazotized sulphanilic acid and an intense color complex was not
formed, while higher than this volume of the base decrease the degree of formation
of the diazonium ion and hence decrease the absorption intensity.

Order of addition

To increase the sensitivity of the present work, the Effect of different orders of
reagent addition were examined. The results obtained that the order of addition:
[hydrochloric acid (A) + sodium nitrite(N) + sulphanilic acid (R) + AMX (S) +
sodium carbonate (B)] gave the more intense color of the azo dye, otherwise a lose
in color intensity take place. This order was used in all subsequent experiments.

86
Effect of solvent of dilution

The Effect of different type of solvents (distilled water, ethanol, and methanol) for
the dilution on the intensity of the color was investigated. It was found that ethanol
(96%) gave higher absorbance than water; but distilled water was selected as the
solvent of dilution due to safety and economic reason.

Effect of temperature

Diazotization and coupling reactions were found to be temperature dependent. Di-

azotization was carried out at 0-5 ?C and coupling reactions was carried out at
room temperature. Above 50 ◦ C there was a decrease in intensity of the color [20].
Therefore, the reaction was examined in cold, room and hot temperature. The re-
sults were indicates nearly the same absorbance when the color developed at room
temperature or when the reaction flask placed in ice bath (2C◦ ) [21]. While at high
temperature (50 C◦ ) the absorbance decrease, for that reason room temperature was
more suitable for this complex formation.

Development time and stability of the azo dye

The Effect of time on the absorbance of the azo dye formed under the optimized
conditions were studied, the solution of the dye stable down to 10 min from the
completion of the volume, after that the color stability of the azo dye gradually
deteriorates. Therefore, the absorbance measurements were carried out after 3 min
from the preparation of the colored azo dye, because the color complex was stable
during 3-10 min.

Final absorption spectra

The final absorption spectrum was recorded after obtaining the optimum conditions.
It was found that the same absorption spectrum (?max = 439.5 nm) was obtained,
as shown in Fig. (2).

Recommended procedure

3.0 ml of AMX solution (50 µgml−1 ) (6.0 µgml−1 AMX in 25 ml) and (2.0 ml)
of (0.25 M) sodium carbonate were added respectively to a 25 ml volumetric flask
containing 1.5 ml of 0.1 M hydrochloric acid, 2.0 ml of was diluted to the mark with
distilled water. The blank solution was prepared in the same way in the absence of
AMX. After 3.0 min the absorbance of the yellow azo dye was measured at 439.5
nm at room temperature.

Calibration graph and statistical data

The calibration graph Fig. (3) obtain by recommended procedure under optimal ex-
perimental conditions described above by plotting the absorbance against the con-

87
 centration of AMX, the colored system obeyed Beer’s law in the concentration
ranges of (0.1 - 11.0 µgml−1 ) of AMX with a detection limit of 0.07 µgml−1 . The
statistical data for the calibration curve such as correlation coefficient, Sandell index
and molar absorptivity were found to be 0.9995, 13.157 µg cm−2 and 2.9442x104
L mol−1 cm−1 respectively.

Figure 3: The calibration graph for the spectrophotometric determination of AMX using diazoti-
zation method

Precision and accuracy

To calculate the precision and accuracy of the proposed method, the standard devia-
tion (SD), relative standard deviation (RSD) and relative error (E%) for ten replicate
of two different standard AMX concentrations were determined. The results were
shown in Table (2) which indicates a good precision and accuracy.

Table 2: Precision and accuracy of the proposed method

Concentration of AMX
Mean absorbance (n=10) SD RSD% E%
(µgml−1 )
3.0 0.2366 6.810x10-4 0.28780 1.13
10.0 0.7658 2.347x10-3 0.30655 - 0.03

Interferences

In order to perform possible analytical application of the present method, Effects

of some foreign species which often accompany with this drug in pharmaceutical
products were studied by adding a known amount of foreign species to 6.0 µgml−1
of AMX. The Effect of interfering species were considered to be significant if
the relative error (E%) less than ±5%. Table (3) shows the maximum tolerance
concentration of the various interfering species; tolerance level in most cases was
≤ ± 5%.

88
 Table 3: Effect of the interferences on the absorbance of 6.0 µgml−1 AMX
Interferences Tolerance level Absorbance E%
Glucose 90 0.441 - 4.90
Galactose 100 0.458 - 2.34
Fructose 100 0.470 0.21
Sucrose 100 0.459 - 2.10
Lactose 100 0.450 - 4.05
Starch 100 0.467 - 0.42
Sodium chloride 100 0.474 1.06

Applications

The present method was applied successfully for determination of AMX in com-
mercial pharmaceutical products. For the aim of comparison the samples were
also analyzed by the standard method in British Pharmacopeia [22] including UV-
spectrophotometric method. Results are shown in Table (4).
The results of two methods were compared using the t-test and F-test with a
confidence limit of 95% indicates that there were no significant differences between
accuracy and precision of the two methods, which indicates that there was no any
systematic error in our method of analysis.
|x| √
Where : tcalculated = n,
SD
x = Mean of difference between two methods.
SD = Standard deviation of difference between two methods.
n = number of sample (degree of freedom).

From the values of tcalculated (0.830) of the present method was less than the ttable
(2.571) [23] with a confidence limit of 95% the results indicated that there is no
significant difference between the accuracy of the two methods.
SD12
Fcalculated = ;
SD22
SD1 and SD2 = Standard deviation of two methods.

The value of F-calculated (1.466) was less than the F-value from the table (5.05)
with a confidence limit of 95%, the results indicated that there is no significant
difference between the precision of two methods.

Structure of the dye

The stoichiometric ratio between AMX and sulphanilic acid in azo dye was checked
by JobŠs continuous variation method [24]. For this method, a series of solutions
was prepared in which the total volume of AMX and sulphanilic acid were kept
at 2.0 ml. Results obtained in Fig. (4) Showed the existence of a (1:1) AMX:
sulphanilic acid reagent at 439.5 nm. Therefore, the structure of the dye may be
written as indicated in Fig. (1).

89
 Table 4: Comparison of present method with standard method [22] for determination of 6.0 µgml−1
AMX using t-test and F-test

Amount found µgml−1 Absolute

Formulation E%
error
Present Standard
Method method [22]
Amoxicillin
S.D.I. 5.724 5.786 - 0.062 - 1.071
Capsule (Iraq)
Glomox
gobalpharma 5.921 5.821 0.100 1.718
(U.A.E.)
Penamox
HIKMA 5.632 5.750 - 0.118 - 2.052
(Jordan)
Amoxicillin
Tablet HEUMANN 5.686 5.607 0.079 1.409
(Germany)
Amoxicillin
ACAI 6.565 6.429 0.136 2.115
Injection
(Iraq)
Largopen
6.105 6.036 0.069 1.143
bilim (Turkey)
Mean 5.938833 5.904833 0.034000
SD 0.353338 0.291816 0.100329
tcalculated =0.830 ttable = 2.571
Fcalculated =1.466 Ftable = 5.050

Figure 4: Job’s plot method of AMX - sulphanilic acid in azo dye

FIA-Spectrophotometric Determination of Amoxicillin Using Diazotized Sul-

phanilic Acid

The FIA system, Fig. (5), with a sample volume of 50 µL used to inject 25.0
µgml−1 of AMX into the diazonium ion solution stream. Four reaction coils were
used in the system with length a = 20 cm, b = 20 cm, c = 10 cm and d = 40
cm. Nitrous acid was produced by the reaction of sodium nitrite with hydrochloric
acid in the first reaction coil, then it was merged with sulphanilic acid to form the
diazonium ion in the second reaction coil, after that the diazonium ion reacted with
AMX (25.0 µgml−1 ) to form the azo dye in the third reaction coil. Finally, this

90
stream merged with sodium carbonate solution in the fourth reaction coil to form
stable yellow color complex in basic medium, with a total flow rate 1.0 ml min−1 .
The absorbance was measured in (mV). At least three injections were made for
every sample solution.

Figure 5: Schematic diagram of the FIA-spectrophotometric manifold used for the determination
of AMX by diazotization method

Optimization of the experimental parameters

Effects of various parameters on the absorbance of the azo dye were studied and
the reaction conditions were optimized. The optimization started using 0.006 M
hydrochloric acid, 0.005 M sodium nitrite, 0.0001 M sulphanilic acid, 0.02 M
sodium carbonate, four reaction coils with length (a = 20 cm, b = 20 cm, c = 10
cm and d = 40 cm), 1.0 ml.min−1 flow rate, 50 µL sample volume and 100 µL
volume of the flow through the cell at the maximum wavelength of 439.5 nm for
the determination of AMX.

Effect of the flow rate

Effect of flow rate was tested in the range of 0.2 - 1.2 ml min−1 in order to obtain
best absorbance intensity of the signals. It illustrates that the peak height decreased
with increasing flow rate, 0.2 ml min−1 gives the higher peak height than 0.5 ml
min−1 , but the flow rate 0.5 ml min−1 was selected as optimum flow rate, because
very low flow rates cause decrease the speed of analysis and increases the time for
the peak to return to the base line and resulted in the formation of a broad peak.
Therefore, 0.5 ml min−1 was used in further work.

Effect of length of reaction coils

The Effect of mixing coil lengths on peak height was studied in the range (0 - 60
cm), four optimum values were obtained. Results show that the optimum condition
was observed for four reaction coils (a, b, c, d) used in the system [Fig.(5)] were
(10, 30, 20, 40 cm) respectively. All reaction coils used in the manifold were made
of Teflon with an i.d. of 0.6 mm.

91
 Effect of sample volume

The Effect of sample volume on the peak height was investigated by injecting
different volumes using 40 - 190 µL of the sample volume. The sensitivity of
the proposed method increased significantly with increasing the sample volume into
the flow system up to 100 µL, after that it decreases with increasing sample volume,
because complete mixing was not occur due to increase the dead volume. Therefore
100 µL sample volume was employed in subsequent measurements.

Effect of hydrochloric acid concentration

The Effect of hydrochloric acid concentration on the peak height was studied in the
range of (0.003 - 0.1 M). It was found that 0.01 M of HCl gave the maximum dia-
zonium ion production. Higher concentration of HCl weakens the alkaline medium,
which required for obtaining the color of the azo dye and at lower HCl concentration
the production of nitrous acid was not enough to form an intense color of the azo
dye. Therefore, this concentration (0.01 M) was selected for further studies.

Effect of NaNO2 concentration

The Effect of sodium nitrite concentration on the color intensity was investigated
by changing the concentration of NaNO2 in the range (0.002 - 0.5 M). The con-
centration of NaNO2 that gave the greatest peak height was found to be 0.05 M.
Therefore this concentration was chosen in further studies. Additional increase of
nitrite concentration produces unstable and noisy signals due to the possibility of
generating bubbles of nitrogen gas.

Effect of sulphanilic acid concentration

The Effect of sulphanilic acid concentration upon the analytical response of the flow
system was studied in the concentration range of (5x10-5 - 1.0x10-2 M). Results
indicate that 5x10-3 M concentration of sulphanilic acid gave maximum response.

Effect of sodium carbonate concentration

Effect of sodium carbonate in the concentration range of 0.001 - 0.5 M was inves-
tigated and a greatest analytical response (peak height) was obtained with 0.01 M.
Low concentrations of the sodium carbonate were not sufficient for producing an
intense color while color intensity decreased at higher concentration, due to increase
the basicity of the solution.

Calibration graph

Using optimum conditions, a straight line of calibration curve was obtained by

plotting the peak height (mV) against the concentration of AMX (µgml−1 ) as shown
in Fig. (6). Beer’s law of the colored system is obeyed over the concentration

92
range of 1.0 - 35.0 µgml−1 of AMX, with correlation coefficient (0.9985) and
the detection limit of 0.7 µgml−1 . The precision and accuracy of the proposed

Figure 6: Calibration graph for the FIA-spectrophotometric Determination of AMX using diazoti-
zation method

method were determined by ten replicate measurements which made on two different
concentrations of the pure AMX solution. The precision of the method was checked
with relative standard deviation (RSD%), while the accuracy was checked with
relative error (E%) of the same solution. The results were shown in Table (5) which
indicates a good precision and accuracy.

Table 5: Precision and Accuracy of the method.

AMX
Mean
concentration SD RSD% E%
(n = 10)
(µgml−1 )
4.0 101.4 0.966 0.950 - 0.95
25.0 624.8 3.119 0.499 - 3.07

Interferences

In order to perform possible analytical applications of the present FIA system,

Effects of possible interferences on the determination of AMX in pharmaceutical
preparations were studied by injecting synthetic sample solutions containing 20.0
µgml−1 AMX. This study was performed by comparing the signal obtained when
AMX present alone and in the present of different interferences. The results obtained
were summarized in Table (6) which shows the maximum tolerable concentration
of the various interferences up to 100 folds greater than the amount of the sample
(AMX) that leading to obtain the tolerance level less than ±5.0%.

Applications

The present FI-spectrophotometric method was applied for the determination of

AMX in commercial pharmaceutical products. For the aim of comparison the

93
 Table 6: Effect of interferences on the peak height of 20.0 µgml−1 AMX
Interferences Tolerance level Peak height (mV) E%
Glucose 100 517.4 - 4.31
Galactose 100 533.3 - 1.37
Fructose 100 530.2 - 1.94
Sucrose 100 520.5 - 3.74
Lactose 100 522.0 - 3.46
Starch 100 527.3 - 2.48
Sodium chloride 100 550.8 +1.87

samples were also analyzed by the standard method in British pharmacopeia [22] as
mentioned in Table (7). The results of the proposed method and standard method
were compared using the t-test and F-test. The student t-test and F-test [23] show that
there is no significant difference between the two methods with regard to accuracy
and precision (tcalculated = 0.025 was less than ttable = 2.571 and Fcalculated =
1.034 less than Ftable = 5.05 with a confidence limit of 95%). Which indicate that
there was no any systematic error in the measurements.

Table 7: Comparison of present method with standard method [22] for determination of 6.0 µgml−1
AMX using t-test and F-test

Amount found µgml−1 Absolute

Formulation E%
error
Present Standard
Method method [22]
Amoxicillin
S.D.I. 18.973 19.143 - 0.170 - 0.888
Capsule (Iraq)
Glomox
gobalpharma 19.648 19.679 - 0.031 - 0.158
(U.A.E.)
Penamox
HIKMA 18.837 19.179 - 0.342 - 1.783
(Jordan)
Amoxicillin
Tablet HEUMANN 18.452 18.179 + 0.273 1.502
(Germany)
Amoxicillin
ACAI 22.849 23.286 - 0.437 - 1.877
Injection
(Iraq)
Largopen
21.824 21.143 + 0.681 3.220
bilim (Turkey)
Mean 20.09717 20.10150 - 0.00433
SD 1.806357 1.836845 0.418562
tcalculated = 0.025 ttable = 2.571
Fcalculated =1.034 Ftable = 5.05

94
Conclusion

In the present work batch and flow injection analytical techniques have been used
for determination of AMX in pharmaceutical formulations. In the batch method
a simple, sensitive, low cost reagent and rapid spectrophotometric method for the
determination of AMX in pharmaceutical formulations (capsules, tablets, and in-
jections) had been devised, it was free from tedious steps like extraction, complex
sample treatment and heating. The present method based on coupling of AMX with
diazotized sulphanillic acid in basic medium. The method was applicable for low
concentration of AMX with competitive precision (RSD%; 0.2878, 0.3065), high
accuracy (E%; 1.13, -0.03) and reasonable sensitivity in which the molar absorptiv-
ity was found to be 2.9442×104 L mol−1 cm−1 and sandell’s index 13.157 µgml−1
for low concentration linear calibration. Beer’s law was obeyed over the range of
AMX concentration, 0.1 - 11.0 µgml−1 .
The sensitivity of this method was not less than that obtained for the determi-
nation of AMX using spectrophotometric methods [5, 25, 26, 27]. These methods
were linear ranges between 2.0 - 20.0 µgml−1 , 5.0 - 13.5 µgml−1 , 5.0 - 30.0 µgml−1
and 5.0 - 60.0 µgml−1 respectively.
FI-spectrophotometric method was successfully used for the determination of
AMX in capsules, tablets, and injections with advantages simplicity, low cost, high
sensitivity, high accuracy and precision. Sensitivity of this method was not less than
that obtained for determination of AMX using FIA methods [6, 8], that gave the
linear ranges between 10.0 - 70.0 µgml−1 and 50.0 - 1200.0 µgml−1 respectively.
As comparison in sensitivity between batch and FIA techniques it seems that
batch technique with linear range of 0.1 - 11.0 µgml−1 and detection limit of 0.07
µgml−1 was more sensitive than FIA technique with linear range of 1.0 - 35.0
µgml−1 and detection limit of 0.7 µgml−1 , while FIA technique gave more precise
results.

Figure 7: Calibration graph for the FIA-spectrophotometric Determination of AMX using diazoti-
zation method

95
 References

[1 ] L. Ohannesian and A. J. Streeter; Chapter five “Handbook of Pharmaceutical

Analysis”, Marcel Dekker, Inc, New York, (2002), 201.
[2 ] British pharmacopeia, 4th Ed., Pharm. Press London, (2009), 353.
[3 ] R. G. Prasad and S. Rao; Journal of Pharmacy Research, 3(4), (2010), 869 -
872.
[4 ] P. Satish and P. Natavarlal; International Research Journal of Pharmacy, 2(9),
(2011), 48 - 51.
[5 ] S. Imam pasha, T. Mohammed, S. kumar and V Nikhila; International Journal
of Pharmacy and Industrial Research, 2(2), (2012), 106 - 109.
[6 ] M. Q. Al-Abachi, H. Hadi and A. M. Al-Abachi; Analytica Chimica Acta,
554, (2005), 184 - 189.
[7 ] M. Q. Al-Abachi and H Hadi; Iraqi Journal of Science, 50(1), (2009), 8 - 15.
[8 ] M. Q. Al-Abachi and S. Subhi; Journal of Al-Nahrain University Science,
16(1), (2013), 42 - 52
[9 ] Y. Li and J. Lu; Luminescence, 21(4), (2006), 251 - 255.
[10 ] A. I. Chivulescu, M. B. Doni, M. C. Cheregi and A. F. Danet; Rev. Roum.
Chim., 56(3), ( 2011), 247 - 254.
[11 ] A. El Walily, A. Gazy, S. Belal and E. Khamis; J Pharm Biomed Anal.,
20(4), (1999), 643 - 653.
[12 ] E. Kilic, F. Koseoglu and M. A. Akayt; Journal of Pharmaceutical and
Biomedical Analysis, 12(3), (1994), 347 - 352.
[13 ] M. M. Issa, R. M. Nejem, M. Al-kholy, N. S. El-Abadla, R. S. Helles and
A. A. Saleh; J. Serb. Chem. Soc., 73(5), (2008), 569 - 576.
[14 ] A. Ashnagar and N. G. Naseri; E-Journal of Chemistry, 4(4), (2007), 536 -
545.
[15 ] A. H. Khuroo, T. Monif, P. R. P. Verma and S. Gurule; Journal of Chromato-
graphic Science, 46, (2008), 854 - 861.
[16 ] A. C. Vispute, V. A. Patil and D. A. Patil; Der Pharmacia Lettre, 5(1), (2013),
327 - 333.
[17 ] K. M. Mahmoud; Ph. D. Thesis, University of Salahaddin-Erbil, (1997), 43.
[18 ] R. O. Hassan and Y. J. Azeez; Tikrit Journal of Pharmaceutical Sciences,
1(2), (2005), 1 - 8.
[19 ] S. Patai; “Supplement F2 The Chemistry of Amino, Nitroso, Nitro and Re-
lated Groups”, John Wiely & Sons, Ltd, (1996), 628.
[20 ] T. Cherian and B. Narayana; J. Braz. Chem. Soc., 17(3), (2006), 577 - 581.
[21 ] M. S. Abdullah; Ph. D. Thesis, Unversity of Salahaddin-Erbil, (2008), 22.

96
[22 ] British Pharmacopia, H. M. Stationary Office, London. (1993), 773.
[23 ] D. C. Harris; “Quantitative Chemical analysis”, 7th Ed., W. H. Freeman and
Company, New York, (2007), 58, 63.
[24 ] C. Y. Huang, R. Zhou, D. C. H. Yang, P. B. Chock; Biophysical Chemistry,
100, (2003), 143 - 149.
[25 ] K. Farhadi, S. Ghadamgahi, R. Maleki and F. S. Asgari; Journal of the
Chinese Chemical Society, 49, (2002), 993 - 997.
[26 ] H. Salem and G. A. Saleh; J. Pharm. Biomed Anal., 28(6), (2002), 1205 -
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[27 ] M. Q. Al-Abachi and H. Hadi; Journal of Al-Nahrain University, 10(2),
(2007), 1 - 6.

97
98
Physics Session
Simulation of Cu Electrodeposition from Aqueous Solution for Solar
Absorption Using COMSOL Model
Azeez Abdullah Azeez Barzinjy1 , ∗, Haidar Jalal Ismae1 , Kadhim Qasim Jabbar1 ,
Mohammed Abdullah Hamad1 and Samir Mustafa Hamad2
1
Department of Physics, College of Education, Salahaddin University-Erbil, Iraq
2
Department of Mathematics, Faculty of Science, Soran University, Soran-Erbil,
Iraq
[email protected]
Abstract
Electroplating plays a significant role in our daily lives. Both current density
and thickness play a vital role in electrodeposition. They depend on various fac-
tors such as composition, temperature, distance between electrodes, geometry of
electrodes etc. Spectrally selected layer is a vital process in solar absorption. For
absorber-reflector cycles, a coating possess high absorptance at solar wavelengths is
deposited onto a highly IR-reflecting metal substrate (such as copper). Additionally,
Cu is the most utilized substrate material for solar absorbers. Therefore this study
focused on copper due to its availability, very low thermal emittance and price.
Electrodeposition of copper can numerically imitated by means of the secondary
electrodeposition module of COMSOL and the outcomes are compared with the
experimental results to verify model accuracy. The current density for the cathode
geometry and the time required for deposition are also assessed. This study explicit
the lab scale studies, validation of simulation and further prediction of current den-
sities for pilot plant scale in aqueous solutions.
Keywords: Electrodeposition, solar absorption, Cu deposition, COMSOL multi-
physics, aqueous solutions.

Introduction

Copper is the topmost common metal plated among the heavy metals [1]. The main
usages of electroplated copper are plating on plastics, printed wiring boards, zinc
die forming, automotive buffers, electrorefining, and electroforming [2-3]. Copper
is an exceptional select for an underplate and it possess very high plating efficiency,
resulting in exceptional treatment even on difficult-to-plate parts [1].
The usage of copper as an undercoat is preferred by the fact that it is a relatively
low-cost metal and also has the very valuable feature that it can be deposited on top
of numerous substrates with slight technological concerns [4].
Aqueous electrolytes, utilized in metal finishing for electrodeposition of metal
coatings, are essentially solutions of metal salts, which, in solution, dissociate to
procedure electrically charged anions and cations. These ions can be simple metal
cations; Mez+ positively charged complex ions or hydrogen ions, H+ . Also existent
will be ionized acid species, negatively charged complexes or hydroxyl anions, OH− .
The electrical properties of a metal are basically linked to its atomic structure.

99
 Extremely conductive metals such as copper, silver and gold, have a single free
electron per atom. Recall that[5] copper would be a respectable substitute to more
expensive gold and silver as a result of its dielectric function analogous with that
of gold in 600Ű750 nm range [6]. Thin film solar cells and batteries have material-
ized as critical requests for lots of procedures described in literatures [7]. Copper-
semiconductor combinations, for instance, possess efficiencies near 14% which is
relatively acceptable in manufactories [7].
This investigation focuses on solar selectivity obtained with optical thin films.
Metals on the whole possess high reflectance that increases with wavelength towards
the infrared according to the Hagen-Rubens relation. This agrees with low thermal
emittance, i.e., the second prerequisite for solar selectivity. Most metals show high
reflectance in the visible range, this will damage the first condition - it will not
be successfully heated. If we choose copper, it signifies a small step in the right
direction. As a result of inter band transitions, copper absorbs some green and blue
light, and has low thermal emittance. Nevertheless, a clean surface of copper does
not become hot enough in sunshine to be of any practical use [5].
Electroless deposition is a reasonably low-cost method as compared to physi-
cal vapor deposition. Additionally, by governing the bath chemistry and working
conditions, extra elements can be combined into the electroless Cu deposits, given
that electroless Cu films enhanced electro-migration resistance and other higher
properties [8-10].
In this investigation a theoretical model for the electroless deposition of copper
on a planar electrode is presented and utilized, by means of COMSOL Multiphysics,
to make time-dependent expectations on the numerous parameters in the system. In
order to verify the model precision the experimental data, described by Ramasub-
ramanian et al. [11], utilized in this investigation as an appreciated study in Cu
electrodeposition from aqueous solution.

Theory

In this study a copper-tartrate-formaldehyde bath is considered for the deposition

process. Copper is entirely complexed by tartrate and exists as Cu(OH)2 L−4 2 ion,
where L is the tartrate ligand. It can be agreed that tartrate is a good complexing
agents spatially in the alkaline system. Hence, Basol et al. carried out two-step
electrodeposition utilizing tartrate and citrate as complexing agents in the alkaline
regime [12]. Therefore, the reactions that were modeled in this investigation are:

Cu(OH)2 L−4
2 + 2e
−
→ Cu + 2OH − + 2L−2 (1)
2HCHO + 4OH − → 2HCOO− + H2 + 2H2 O + 2e− , (2)

The geometry considered in this model is consist of the electrode surface at one
end, diffusion layer of thickness 2 mm and the bulk electrolyte at the other end.
Fig. 1 demonstrates the representation diagram of the electroless deposition pro-
cedure proposed in this investigation. For the experimental settings [11] and the
reactions stated above, one could assume the significant reliant on variables that

100
Figure 1: A schematic representation of the electroless plating modeling region on planar electrodes.

must be followed in the system. These include eight types of ions; Cu(OH)2 L−4 2 ,
HCHO, HCOO− , H+ , OH− , SO−2 4 , Na +
and L −2
,[11] the solution potential (Φ),
and the electrode potential (V). The transport of 8 ionic species through diffusion
is explained in the model.
In the diffusion layer, the concentration of each ionic type is ruled by the follow-
ing material balance equation:
∂Ci
= −∇.Ni + Ri , (3)
∂t
where Ri is the degree of standardized production of species i from entire reaction.
The flux of each species i is specified by:
zi Di F Ci ∂Φ ∂Ci
Ni = − − Di . (4)
RT ∂y ∂y
This equation comprises only the influences of migration and diffusion. In solution,
water stability exists at all times and therefore the equation:
CH + COH + − Kw = 0, (5)
must be fulfilled. The potential in the solution changed in agreement with the
electro-neutrality circumstance set by:
X
zi Ci = 0. (6)
At the diffusion layer-bulk interface, y = ∞, the concentrations of every single
species of the seven species relates to its bulk values:
Ci = Cib (7)
The bulk equilibrium conditions are governed by the equations described by Rama-
subramanian et al. [11]. At the electrode surface, y = 0, the flux for the each of
the ionic species can be written as:
sij ij
Ni = − , (8)
nj F

101
 where, ij is the current density of whichever the reduction or oxidation reaction.
The cathodic current is supposed to have a negative sign, and the anodic current is
supposed to have a positive sign. sij is the stoichiometric coefficient, balance the
number on both the reactant and product sides, of species i in the electrochemical
reaction j which is stated for the reduction reaction, equation (1) as:
X
sij Mizi → nj e− . (9)
i

In the case of the formaldehyde oxidation reaction, equation (2), the stoichiometric
coefficient is measured by writing the reaction in a divergent way, as:
X
nj e− → sij Mizi . (10)
i

The whole copper reduction reaction, equation (1), and the formaldehyde oxidation
reaction, equation (2), may each contain two or more straightforward steps, the exact
kinetic form of which is not known with certainty. Nevertheless, one can combine
the currents due to each of these elementary steps, as outlined by Newman [13], and
achieve expressions for the whole reaction rate with regard to the concentrations of
the electro-active species. The partial current densities of the cathodic and anodic
reactions can be expected to follow Butler-Volmer kinetics. If the reduction reaction
goes through a number of steps, all of which are in series, and the intermediate
species formed do not diffuse away from the electrode, the partial current density
for copper reduction can be written as:
(  Pi1    qi1
Ci,o αa1 F  θ
 Ci,o
i1 = io1,ref Πi exp V − Φo − U1 − Πi
Ci,ref RT Ci,ref
 
−αc1 F 
V − Φo − U1θ

exp , (11)
RT
where pi1 = si1 when si1 > 0 and qi1 = si1 when si1 < 0. The relations in
parentheses in the equation (11) indicate the driving force for a one electron transfer
reaction to the Cu(OH)2 L−4 2 complex forming an intermediate. By means of the
above-mentioned settings of zero intermediate loss and sequence reactions, the total
current density for the cathodic reaction can be measured by multiplying a factor of
two, which represent number of electrons involved in the reaction, to the fundamental
Butler-Volmer relation in equation (11) [13]. It is also worth noticing here that the
concentration profiles of the numerous components rely only upon the total reaction
rate and therefore the precise kinetics of these reactions is not essential for this
determination. A comparable dealing of the total formaldehyde oxidation reaction
results in the current density for the formaldehyde oxidation reaction to be given
by:
(  Pi2    qi2
Ci,o αa2 F  θ
 Ci,o
i2 = io2,ref Πi exp V − Φo − U2 − Πi
Ci,ref RT Ci,ref
 
−αc2 F 
V − Φo − U2θ

exp , (12)
RT

102
where pi2 = si2 when sil < 0 and qi2 = si2 when sil > 0.
In equation (11) and (12) the driving force for the electrochemical reaction is given
by the potential difference term[V − Φo − U1θ ]. V is the electrode potential, Φo is
the solution potential together with the electrode surface, and U1θ is the open-circuit
potential for the reaction i. The open-circuit potential is measured for these two
reactions by a Nernst equation as:
RT Cu(OH)2 L−4 2
U1θ = U1o + ln (13)
nF [OH − ]2 [L−2 ]2
RT [HCOO− ]2
U2θ = U2o + ln (14)
nF [OH − ]4 [HCHO]2
where U1o and U2o are the typical electrode potentials for the reduction and oxidation
reactions, in that order. Meanwhile, throughout electroless plating, there is no net
current flowing through any outward system, the cathodic and anodic currents must
be identical and opposite to each other and therefore the resulting equation has to
be fulfilled at all times.

i1 + i2 = 0 (15)

The least thickness, whereupon there is a minor jump in the concentration profile
between the bulk-diffusion layer interface and the nearby node points in the diffusion
layer, was expected to be the thickness of the diffusion layer. This thickness was
characteristically about 0.2 cm for all the conditions demonstrated in this investiga-
tion. The electrode potentials supposed for the copper reduction and formaldehyde
oxidation reactions are 20.266 and 1.5 V, correspondingly.

Model Parameters

The conveyance and kinetic parameters necessary for implementation these inten-
tions are in depth in Table 1. The bulk concentrations of each variable under an
assumed set of settings were attained by means of the process detailed in a preceding
investigation [15].

Table 1: Physical data and model parameters used in this investigation

Physical data
Species
zi Di × 10 (cm /s) si,1 si,2 Ci,bulk (mol/cm3 )
5 2

Cu(OH)2 L−4
2 -4 0.7a -1 0 0.06 × 10−3
HCHO 0 1.2a 0 -2 0.22 × 10−3
L−2 -2 0.794b 2 0 0.39 × 10−3
SO−2
4 -2 1.065b 0 0 0.06 × 10−3
OH− -1 5.273b 2 -4 0.16 × 10−4
a
From Ref. [11].
b
From Ref. [14]

103
 Results and Discussion

Fig. 2 demonstrates the expected current density plot for the electroless copper
plating from a bath containing 0.1 M CuSO4 , 0.2233 M HCHO at a pH of 12.25.
It can be seen that, the current density is noticeably high in the early stages of elec-

Figure 2: Electroless copper deposition current density as a function of time for a bath containing
0.1 M CuSO4 , 0.2233 M HCHO at a pH of 12.25.

troless deposition process, which is attributed to greater driving force. The partial
current density plot indicates a very important reduction in current density in the
early time periods carefully causing the potential plot presented in Fig. 3. The al-
teration in the electrode surface potential is credited to the equilibrium potentials of
copper reduction and formaldehyde oxidation reactions. At very short periods, the

Figure 3: Predicted electrode potential values during electroless plating for various concentrations
of copper and formaldehyde in the bath.

driving force for the reaction, the modification between the open-circuit potential
and the mixed potential for the cathodic and anodic reactions, is considerably high.

104
Therefore, the greater driving force causes a large current. As a result of the move-
ment of the open-circuit potentials in the positive direction and the concentration
reduction at the surface, the polarization decreases with time producing the current
to reduce promptly. At following times, the electrode potential become stable, this
again causes the current and the corresponding surface concentrations to reduce.
Fig. 4 illustrates a plot of the concentrations of three species at the surface for
limited plating times. The concentrations of these three species, which get con-
sumed at the surface, drop dramatically throughout the initial periods. Then the
concentrations increase as a result of the influx of new ions from the bulk by diffu-
sion and migration. The preliminary reduction in the concentrations of Cu (OH)2
L−4
2 and HCHO is initiated by the fast preliminary reaction rate. This reduction
is counterbalance by the diffusion and migration components of these species to a
small range, which causes the surface concentrations to recuperate. Likewise, the
concentration of free tartrate increases suddenly and afterward becomes constant
because of diffusion toward the bulk. The normalized concentrations of Cu (OH)2

Figure 4: Concentrations of the copper complex, formaldehyde, and the free tartrate ligand (L−2 )
at the electrode surface versus time.

L−4
2 , HCHO, and L
−2
can be divided into three regions: (i) the preliminary region
from 10 to 100 s where the surface concentrations decrease for Cu (OH)2 L−4 2 and
−2
HCHO and increase for L ; (ii) the intermediate time zone region between 100
and 500 s where the concentration levels off; (iii) the longer time regions where the
concentration increase slightly in L−2 and decreases significantly in Cu (OH)2 L−4 2
and HCHO cases.
The copper deposit thickness as expected by the COMSOL model is plotted for
the solution at a pH of 12.25 as indicated in Fig. 5. In the mentioned pH value
the model was run for, the partial current density was very high in the beginning,
and consequently reduced considerably as observed by the variation in slope of the
plating thickness. Schoenberg indicated that the maximum rate of deposition for
the tartrate solution is at pH 12.8 [16]. Later the current become stable and a linear
relationship of the plating thickness was obtained with time. Experimental investi-
gations carried out in the literatures [16, 17] have shown that there is a maximum

105
 that exists for the plating rate versus a pH values. This has been recognized to
the separation of methylene glycol, which is made under circumstances of high pH
values.

Figure 5: Electroless copper deposit thickness as a function of time at a pH of 12.25.

Conclusions

This theoretical investigation focused on Cu since Cu is the most utilized substrate

material for solar absorbers. Solutions comprising tartrates may be functioned with
lower free cyanide and at higher current densities and efficiencies without damaging
anode corrosion. The model expects by means of COMSOL Multiphysics, the
surface concentrations, electrode potential, and the plating rate as a function of
different electro-active species concentrations in the plating bath.
The concentration of species such as copper complex and formaldehyde decreases
with time because they are used up throughout the deposition whereas that of the
free tartrate ligand increases with time because it is created in the deposition pro-
cess. The expectations achieved for the electroless copper deposition from tartrate
baths indicate that OH− ion concentration is limited by mass-transfer and therefore
dominates the plating rate. The obtained results in this study, in general, are in a
good agreement with available experimental dada.

Acknowledgment

The authors would like to thank Dr Jamil Juma at Materials Centre/ University of
Leicester, Leicester-UK, for his insightful ideas, suggestions, and time for serving in
our discussions. If anyone who needs the details about the COMSOL model please
contact the corresponding author for any additional help.

106
References

[1 ] Flott, L., The electroplating process. Metal Finishing (USA), 1996. 94(3):
p. 55-58.
[2 ] Schlesinger, M. and M. Paunovic, Modern Electroplating. 2011: Wiley.
[3 ] Garrou, P., C. Bower, and P. Ramm, Handbook of 3D Integration: Volume
1 - Technology and Applications of 3D Integrated Circuits. 2011: Wiley.
[4 ] Gamburg, Y.D. and G. Zangari, Theory and Practice of Metal Electrodepo-
sition. 2011: Springer New York.
[5 ] Piegari, A. and F. Flory, Optical Thin Films and Coatings: From Materials
to Applications. 2013: Elsevier Science.
[6 ] Chan, G.H., et al., Plasmonic properties of copper nanoparticles fabricated
by nanosphere lithography. Nano Letters, 2007. 7(7): p. 1947-1952.
[7 ] Martin, P.M., Handbook of deposition technologies for films and coatings:
science, applications and technology. 2009: William Andrew.
[8 ] Inoue, F., et al., Electroless Cu deposition on atomic layer deposited Ru as
novel seed formation process in through-Si vias. Electrochimica Acta, 2013.
100: p. 203-211.
[9 ] Yu, L., Electroless Deposition of Copper and Copper-Manganese Alloy for
Application in Interconnect Metallization. 2014, Case Western Reserve Uni-
versity.
[10 ] Tracton, A.A., Coatings Technology: Fundamentals, Testing, and Processing
Techniques. 2006: CRC Press.
[11 ] Ramasubramanian, M., et al., Solution equilibrium characteristics of electro-
less copper deposition on thermally-activated palladium-catalysed polyimide
substrates. Journal of applied electrochemistry, 1998. 28(7): p. 737-743.
[12 ] Aksu, S., J. Wang, and B.M. Basol, Electrodeposition of InŰSe and GaŰSe
thin films for preparation of CIGS solar cells. Electrochemical and Solid-State
Letters, 2009. 12(5): p. D33-D35.
[13 ] Newman, J., Electrochemical Systems, 1991. Prentice-Hall, Englewood
Cliffs, New Jersey.
[14 ] Lide, D.R. and H.V. Kehiaian, CRC handbook of thermophysical and ther-
mochemical data. Vol. 1. 1994: CRC Press.
[15 ] Shacham-Diamand, Y., V. Dubin, and M. Angyal, Electroless copper depo-
sition for ULSI. Thin Solid Films, 1995. 262(1): p. 93-103.
[16 ] Schoenberg, L., The structure of the complexed copper species in electro-
less copper plating solutions. Journal of the Electrochemical Society, 1971.
118(10): p. 1571-1576.
[17 ] Duffy, J., L. Pearson, and M. Paunovic, The effect of pH on electroless copper
deposition. Journal of The Electrochemical Society, 1983.130(4): p. 876-880.

107
108
 Neutronic and Thermal - Hydraulic Calculations for The Conversion of
BN-350
Jamal Mohammed Rashid Abda and Abdul Karim Thamir Mohammed2
1
1
Department of Physics, Faculty of science and Education sciences, University of
Sulaimani, 46001 Sulaimani, KurdistanŰIraq
2
Baghdad university, Baghdad, Iraq
[email protected]
Abstract
This paper represents the neutronic and thermal-hydraulic calculations for the
conversion of BN-350 from 27% to 7% enriched Uranium fuel elements. Mixed
enrichment cores has been studied where low enrichment uranium (LEU) cores fresh
fuel elements substitute gradually the high enrichment uranium (HEU) depleted fuel
elements in the equilibrium core. Thermal hydraulic calculations have been carried
out to determine changes in the characteristics of the converted reactor during steady
state conditions and transient response to a coolant flow loss (LOCA).
Keywords: neutronic and thermal-hydraulic, BN-350, low enrichment uranium
(LEU), coolant flow loss.
Abbreviations:
HEU : High enriched uranium LEU : Low enriched uranium
DAIXY : A programming code JAM1 : Neutronic computer software
JAM2 : Thermal hydraulic program Kef f : Multiplication factor
` : prompt neutron lifetime βef f : effective delayed neutron fraction
Λ : neutron generation time

Introduction

Several studies for the conversion of BN-350 loop type reactor has been studied.
In the case of this purpose, it was established[1-3] that, in order to match the cycle
length of the current 27% high enriched uranium fuel design HEU with 20%en-
riched fuel LEU( low enriched uranium ), an uranium density of about 11gm/cm3
is required.

The BN-350 reactor will be converted to use the 7% enriched uranium LEU
fuel elements which have only minor changes of the fuel plates and no change
in the design of the fuel element geometry. With the same element geometry,
the thermal hydraulic characteristics of the core should be almost identical with
both HEU and LEU fuels. Neutronic calculations have been performed using well
verified computer codes, DAIXY[4,5] another program established ( JAM1 ) using
a personal computer.

A thermal hydraulic program ( JAM2 )[6] has been developed to analyze the
behavior of the reactor in steady conditions at nominal pond during primary pumps
failure. The following table, table(1) gives some important physical properties of
the BN-350 reactor[7-9]:

109
 Table 1: Reactor and fuel element design descriptions with the HEU and LEU fuels.
Property HEU LEU
Reactor type Loop type Loop type
Steady Űstate power level, Mw
1000 1000
U235 density in fuel meat:
3
External plate, g/cm 8 8.11
Internal plate, g/cm3 9.5 9.523
Uranium enrichment, % 27 20
Number of reference core configuration 211 211
Control blade material B4 C B4 C
Number of control blades 12 12
Number of fuel assemblies 120 109
Coolant Na Na
Reflector Depleted UO2 Depleted UO2
U235/fuel element, kg 1170 1185
Lattice pitch, mm 98 98
Primary coolant velocity, m/s 8 8
Inlet temperature, C◦ 0 300 300
Outlet temperature, C◦ 0 500 500
Meat thickness, mm 6.1 6.1
Meat width, mm 1060 1182
Active length, mm 1060 1060
Clad thickness SS316, mm 0.35 0.35
2
Coolant flow area, cm 5450 5450
Fuel meat composition UO2 PuO2 UO2 PuO2

Neutronic calculations

Cross sections were prepared for different regions in the core using ANL-5800[10,11],
reactor physics constants. Different cell models were needed to generate appropriate
cross sections for the various reactor regions in the standard four group structure.
Most of the results of this study are based on two dimension multigroup diffusion
calculations using DAIXY code and JAM1 personal program, with bucklingŠs im-
posed for the axial dimension [12]. Fig.(1), shows the core arrangement of BN-350
reactor consisting of three regions of 898 fuel elements with depleted UO2 reflected,
the experimental value of Kef f estimated for the reference HEU core using UO2 PuO2
was 1.0288065.

110
 The DAIXY and JAM1 were used to calculate the Kef f for both HEU and LEU
reference core with 898 fuel elements using the same axial buckling[3]. The Kef f
results are shown in table(2) for LEU and HEU cases, in which the code predicts
well the Kef f experimentally determined.

Table 2: Effective multiplication factor calculated for BN-350 with HEU and LEU fuels.
Daixy JAM1 References[5,7,8,12]
Case
Kef f %∆k/k Kef f %∆k/k Kef f %∆k/k
HEU 1.0321169 2.883400 1.0329612 2.0806 1.028806 2.8
LEU 1.0301165 0.117022 1.0309513 0.9040 1.010101 0.9

The reactor kinetics parameters have been calculated for the reference core and at
the end of equilibrium (3) cycle for both HEU and LEU fuels [7,13,14]. The neutron
generation time (Λ). The prompt neutron lifetime (`) and the effective delayed
neutron fraction (βef f ), using the two dimensional diffusion theory perturbation
capability of JAMI, are given in table 3.

Table 3: Kinetics parameters for the HEU and LEU cores.

Core Fuel type Λ(µs) `(µs) βef f
HEU 42.0 43.2098 0.0158
Reference
LEU 39.0 39.3939 0.0146
HEU 43.5 43.7618 0.0149
Calculated
LEU 40.3 40.3472 0.0141

The isothermal temperature and void coefficients of reactivity were computed

separately as functions of temperature or void fraction due to the following effects:
1. Hardening of the neutron spectrum caused by increasing the temperature of
the coolant only.
2. Increasing in neutron leakage when the coolant density is decreased., and
3. Increasing in U-238 epithermal resonances absorption due to the increasing of
the fuel meat temperature ( Doppler effect )[15].
The global temperature coefficient of reactivity, expressed in terms of −∆ρ/∆T ∗
10−3 /C ◦ , obtain for the reference fresh fuel core was 1.761 for HEU fuel and 1.532
for the LEU fuel. The control worth calculations of the 12 fully inserted B4C-SS316
control blades for the reference and equilibrium , both HEU and LEU fuels , are
based on a diffusion calculation. Cross sections were used from ANL-5800 and
from[10,11]. The reactivity ρ corresponding to the control blades using DAIXY
code, are shown in table(4) for both reference and equilibrium cores using HEU
and LEU fuels. Mixed enrichment cores might be considered as an option for the
conversion of the BN-350 reactor where LEU fresh elements substitute gradually
the HEU fuel in the equilibrium core. This situation is expected to occur during the
core conversion and in the planning of such conversion it has to be able to predict
the accurate behavior of the mixed cores.
For the mixed core calculations it will be used the same shuffling pattern that
was utilized for the equilibrium core.

111
 Table 4: Control worth of the B4 C control blades for HEU and LEU cores.
Core Fuel type K∗out K∗∗
in Control worth
HEU 1.0288065 0.9153320 0.1134745
Reference
LEU 1.0309132 0.8631890 0.1677242
HEU 1.0321169 0.9831417 0.0489752
Calculated
LEU 1.0302966 0.9528554 0.0998154
∗
Kout : Kef f when control rods up.
∗∗
Kin : Kef f when control rods down.

Thermal Ű hydraulic calculations

The conversion of the BN-350 reactor from HEU fuel to LEU fuels was considered
without changes in the fuel element geometry. Only minor modifications in fuel
plates were made where the width of the meat was changed from its actual values
of 6.1mm to 7.256mm in the LEU case. Thus, the thermal-hydraulic behavior of
the converted core would be virtually identical to the HEU core [16].
The thermal-hydraulic calculations have been carried out using a thermal-hydraulic
subroutine program JAM2 in order to determine changes in the characteristic of the
converted reactor during steady conditions and transient response to a loss of coolant
flow ( LOCA ) for both HEU and LEU reference core. However, the calculations for
the reference core, both HEU and LEU fuels, confirm that there are minor changes
in the thermal-hydraulic behavior, table(1), and for this reason it will only be shown
the results for the fresh LEU reference core.
For steady state case, the results of temperature distribution at the surface of fuel
rod and the centre were 670C◦ and 369C◦ respectively as shown in fig(2), which is
indicates that temperature was fixed after about 10 seconds. But for transient case,

Figure 1: Temperature distribution with time for fuel rods.( Steady state ).

we show that after 151s that the fuel fused because it reaches its melting point, and

112
the clad fused after 62s for the same reason, as shown graphically in fig.(3).

Figure 2: Transient Temperature distribution in fuel rods ( Fuel and Clad ).

Conclusions

In every aspect, the BN-350 reactor can be converted to use 7% enriched uranium
without modification in the design of the fuel element geometry where the fuel
density changed from its actual values of 9.5gm/cm3 to 9.523 gm/cm3 . The most
important neutronic effect in the equilibrium core performance as a result of the
conversion, from 27% to 7% enriched uranium fuel, is the increasing of thermal
flux from 10.1653×1015n/cm2 .s , to 10.43668×1015n/cm2 .s in the irradiation po-
sitions, and from 7.721287×1015n/cm2 .s , to 8.156372×1015n/cm2 .s for control
rods. Mixed enrichment cores might be considered as an option the conversion of
the reactor where LEU fuel elements substitute gradually the HEU fuels in the equi-
librium core. The thermal-hydraulic behavior for the reference core was identical to
both HEU and LEU fuels. The effects of flow transient over the thermal-hydraulic
characteristics of the reactor for the LEU reference core at beginning of life have
demonstrated that it can operated at 1000Mwt. For the equilibrium core, the flow
through the fuel element channels has a reduction estimated in 10% in a relation
with the reference core, consequently, the maximal heat flux is reduced in about
60% comparing the same configurations. It could be, thus, deduced that the reactor
using equilibrium core is safer than the reference cores configuration, however, this
conclusion must be confirmed.

Acknowledgments

It is a pleasure to thank Dr. Hamid Al-Bahily for the idea of this article developed
from initial discussions with him and important discussions about all difficulties

113
 through preparing this research.

References

[1 ] K.Wirtz, “Lectures on Fast Reactors”, Karlsruhe,(1978).

[2 ] Alan E.Walter and Albert B.Reynolds, “Fast Breeder Reactor”, New York
Pergamon Press, (1981).
[3 ] IAEA, International Atomic Energy Agency “Status of liquid metal cooled
fast reactor technology”, IAEA- TECDCC-1083, (1999).
[4 ] Chantal Vasseur, Robert Bruyere, etal. “Specification of a Diffusion Code in
X-Y Geometry: DAIXY”, CEA Report R2759, (1965).
[5 ] Y. Oka , “Nuclear Reactor Design, An Advanced Course in Nuclear Engi-
neering 2”, DOI 10.1007/978-4-431- 54898-0 2, ľ Authors (2014).
[6 ] SaŠeed, Ghazi Kamal, “Temperature Distribution Calculation in Power Re-
actor Fuel Rod”, Baghdad University, College of Science, (1986).
[7 ] IAEA, “Status of Liquid Metal Cooled Fast Reactor Technology”, INTERNA-
TIONAL ATOMIC ENERGY AGENCY, Vienna, IAEA-TECDCC-1083,ISSN
1011-4289,april, (1999).
[8 ] IAEA, International Atomic Energy Agency “Liquid Metal Cooled Reactors:
Experience in Design and Operation”, IAEA-TECDOC-1569, (2007).
[9 ] S.G.Popov, J.J. Carbajo, V.K. Ivanov and G.L. Yoder, “Thermophysical Prop-
erties of MOX and UO2 Fuels Includind the Effect of Irradiation”, Oak Ridge
National Laboratory, ORNL/TM-2000/351,(2000).
[10 ] ANL-5800, “Reactor Physics Constants”, Argonne National Lab., 2nd ed.,
USAEC, (1963).
[11 ] R.B.Pond and J.E.Matos, ANL/RERTR/TM-26, Argonne National Labora-
tory, Argonne, IL 60439-4841 USA, (1996).
[12 ] IAEA, “Guidebook on Research Reactor Core Conversion From The Use of
High Enriched Uranium to the Use of Low Enriched Uranium Fuels”, Report
IAEA- TECDOC-233,A980.
[13 ] M.A. Shultz, “Control of Nuclear Reactors and Power Plants”, McGrow-Hill
Book Company. Inc., (1961).
[14 ] IAEA BULLETIN, “Liquid Metal Fast Breeder Reactors”, VOL 20, NO
6,(1964).
[15 ] J. Klein, etal, “Neutronic and Thermal-Hydraulic Calculations for the Con-
version of 5Mw LA Reina Reactor Using MEU Fuel”, Childen Nuclear Energy
Commission, Nva- Bilbao 1250, (1985).
[16 ] E.D.Waters , “Heat Transfer Experiments for Advanced Test Reactor”, BWNL-
216, ( 1966).

114
 Isolation, identification and determination of Amygdalin from the sweet and
bitter almond kernel fruit in Erbil city
Kamal M. Mahmoud* and Roza T. Yasin
Salahaddin University- Erbil (SUE), College of Science, Department of Chemistry,
Kurdistan Region, Iraq
[email protected]
Abstract
This paper includes the isolation, identification and determination of amygdalin
from the sweet and bitter almond kernel fruit of Erbil city. Amygdalin has been iso-
lated from the almond kernel fruit using ultrasonic bath, soxhlet and binary solvent
extractions, the effect of temperature and pH has been investigated on the quantity
of extracted amygdalin. TLC, FTIR and HPLC techniques were used to identify
amygdalin qualitatively, while standard pure amygdalin was used as reference com-
pound for the comparison.
Keywords: Isolation, Identification, Almond, Amygdalin, Ultrasonic, Soxhlet, Bi-
nary solvent.

Introduction

Plants contain many active compounds such as alkaloids, steroids, tannins, glyco-
sides, volatile oils, fixed oils, resins, phenols and flavonoids which are deposited in
their specific parts such as leaves, flowers, bark, seeds, fruits, root, etc. [1].
Glycosides are compounds that they chemically consist of a sugar portion (moi-
ety) attached by a special bond to one or more non-sugar portions. Glycosides can
be considered as sugars that are capable of forming ethers with other alcohols, or
esters with acids. Glycosides are broken down upon hydrolysis with enzymes or
acids consisting of non-sugar and a compound sugar in the same molecule: a sugar
moiety (glycone); and non-sugar moiety (aglycone/active portion) [2].
Amygdalin is a cyanogenic diglucoside (D-mandelonitrile-β-D-gentiobioside).
It has the following empirical formula (C20 H27 NO11 ) [3], which has diglucoside
combine with the cyanogenic group, while monoglucoside (prunasin) present in
vegetative organs [4]. Fig. (1) Shows the chemical structure of amygdalin. It

Figure 1: Chemical structure of amygdalin.

115
 was found that the cyanogenic glycoside (amygdalin) present in seeds and leaves of
Elateriospermum tapos. Amygdalin was detected in the samples (seeds and leaves)
using isocratic HPLC with C18 column (150 mm × 3.9 mm), flow rate of 0.8 ml
min−1 , using acetonitrile: water (20:80) as an eluents with 20 µl of sample, and
detection under UV at 218 nm. The amygdalin peak appeared at a retention time of
(2.2 min). The highest amygdalin content was observed in the fresh seeds, (660 µg
ml−1 ), in the cooked seeds, (100 µg ml−1 ), in the fresh leaves (30 µg ml−1 ), and
similar to the amount found in the fermented seeds, (25 µg ml−1 ) [5].
Cyanogenic glycosides are a large group of secondary metabolites that are dis-
tributed in some of plant seeds, kernel, and food products in Kurdistan region,
which are commonly consumed by humans. Amygdalin is one of the cyanogenic
glycosides found, for example, in apples, apricot, and almonds. Therefore, it is
necessary to find sensitive, economic, accurate, rapid and precise analytical meth-
ods for amygdalin qualitatively to investigate extraction efficiency and to determine
levels in a range commercially-available food for the first time. This study was
established an optimal method for the extraction of amygdalin from almonds, iden-
tification of amygdalin from Iraqi-Kurdistan region (Erbil) sweet and bitter almond
kernel (Prunus dulcis).

Experimental

Chemicals, Solvents and Reagents

The chemicals that were used for the detections were of analytical grade. Deionized
water was used throughout. The standard amygdalin was purchased from (Germany).
The chemicals that were used in this study and their suppliers were summarized in
Table (1).

Table 1: Solvents, chemicals and reagents used

Solvents & chemicals Purity % Company
Diethyl ether 99% Scharlau
Ethanol 99% 99% Scharlau
Methanol 99% 99% Scharlau
Acetone 99% Merk
Hydrochloric acid 36-37% Merk
Acetonitrile 99% Scharlau
Potassium chloride ... Merk
Sodium chloride ... Merk
Nitric acid 65% Merk

Apparatuses and Instruments

Soxhelt apparatus with heating Electro thermal EMO 500/30 MK1, (ENGLAND).
IR-Spectrometry: Pie Unicom SP-3005 (USA) spectrophotometer with the range
600-4000 cm−1 was used for the recording IR spectra. Rotary evaporator: laboratory

116
4000, Heidolph company (UK). Ultrasonic bath: Type (TPC-25). High-performance
liquid chromatography HPLC type (Agilent, German) equipped with UV-Vis detec-
tor DAD model SPD-10A(C), Degasser (DG-1310), injection valve containing 20-µl
and 50-µl sample loop. The system is equipped with a commercial column type
LC18 (250 mm × 4.6 mm, 5.0 µm) model (PRONTOSIL, Germany). HPLC-MS
spectrometry (TRACEMS2000, FINNIGAN, USA) was used for differentiation with
HPLC/UV-DAD detector. Hot plate magnetic stirrer: Galleukamp-400 hot plate
magnetic stirrer (UK). Oven: Kottermann ¡ 300 ◦ C (USA). Centrifuge: (Kokusan
H-19F, Japan).

Sampling

The almond kernel food (with both types (sweet and bitter)) was taken from Kurdis-
tan Region in Hawler-market, valley of Balesan, between November and December
2013. The samples were stored in air-tight containers protected from light at room
temperature. Dried plant material (almond kernel) was grounded in a beaker to
obtain a fine powder.

Extraction and Fractionation of the Almond

Ultrasonic Bath System

From (un-pilled sweet and bitter) Plant material (15.0 g) of each samples were
immersed in 150 ml of solvents (using different types of solvent such as diethyl
ether for defatting, ethanol and methanol respectively according to their different
polarity), 0.1 g of active carbon (Norit CNR 115) was added to the solute, and put
it in an ultrasonic bath for 30 min at room temperature, then soaking for 30 min
to reach the equilibrium between solvent and extracting. Extraction was repeated
three times. The extract was obtained after filtration using Buchner funnel [6].
The plant extracts were evaporated in rotary evaporator under reduced pressure to
remove solvent [7,8].

Soxhlet Extraction

A powder of (18.0 g) of each un-shells sweet and bitter almond kernels were ex-
tracted three times with 250 ml methanol for 5.0 h with 0.1 g of active carbon
(Norit CNR 115) using soxhlet apparatus. The mixture was heated by electromag-
net heating. The extracted solvent was obtained after filtration [6] through a filter
paper (Whatman 42). The filtrate was concentrated using rotary evaporator.

Binary Solvent Extraction

Binary solvents were used for extraction of un-pilled sweet and bitter almond kernel.
In this study 1.0 g of the samples were put in a small beaker (250 ml) then added 10
ml of methanol-water mixture in a different ratio [(10:0), (8:2), (6:4), (5:5), (4:6),
(2:8), (0:10)] respectively, the contents of the beaker were shacked continuously

117
 using magnetic stirrer for increasing the extraction process and efficiency [9] for
(30 min) at room temperature. Finally add 0.1 g of active carbon (Norit CNR
115) to the extracted samples, to ensure the freedom from pigments, which may be
interface with the chromatography [10].

Qualitative Analysis of Extracted Amygdalin

Thin Layer Chromatography (TLC)

TLC was used for the identification of amygdalin in the almond. The stationary
phase which was an adsorbent made of silica gel. Reference solution was prepared
according to European pharmacopoeia [11]; 10 mg from amygdalin was dissolved in
2-3 ml water and diluted to 10 ml with methanol. Sample solution was also prepared,
then they spotted onto the TLC plate and placed in a chamber with eluting fluid
(methanol: water (50:50)) as the best mobile phase (after pretested for different
ratios). The eluting fluid, or (mobile phase) rises up the plate via capillary action,
carrying the analytes with it. After ascending development, staining and drying,
silica gel plates were scanned by UV detector at 254 nm. Table (2) shows the Rf
values of standard and sample.

Table 2: Rf values of amygdalin

Compound Rf Color with UV-lamp at 254 (nm)
Reference: Amygdalin 0.53 Green
(Sweet and Bitter) Sample 0.51 Green

Quantitative Analysis of Extracted Amygdalin

Effect of Temperature on the Quantity of Extracted Amygdalin

Extraction of amygdalin was performed by mixing the sample (un-pilled) sweet

and bitter almond kernel with quantity of water (1:10) for a fixed period of time.
A heated magnetic stirring plate was used to preheat the water to 25, 50, 90◦ C,
for stirring the suspension and maintaining the temperature throughout the extrac-
tion. However, the common techniques for extracting amygdalin from foods for
the analysis of amygdalin in the literature are ultrasonic or Soxhlet extraction with
methanol, this is because some amygdalin was decomposed into Benz aldehyde,
HCN and glucose by emulsion a hydrolysis enzyme, and some are converted into
its epimers, neoamygdalin (L-mandelonitrile-β-D-gentiobioside), during the process
of decoction in water [6]. Therefore, in this study the extraction of amygdalin was
performed by three procedures: ultrasonic, soxhlet and extraction by solvent at room
temperature.

Effect of Time on the Quantity of Extracted Amygdalin

The pH and temperature were stabilized for studying the effect of duration of ex-
traction (ultrasonic 30 min, soxhelt 3.0 h and binary solvent extraction 15 min). The

118
main factor that influenced the amygdalin extraction yield is extraction times a short
extraction time would decrease the extraction yield of amygdalin. A long extraction
time would cause some amygdalin to convert into its epimers or decompose in the
aqueous solution [6] and also consuming much more solvents. In this study the
extraction has been repeated three times to make sure that the analyte is achieved
and it is much better than extracting once but for a long time (e.g. more than 3.0
h).

Effect of Solvents on the Quantity of Extracted Amygdalin

A number of different solvents were tested, by changing the volume ratio of the
solvent, to obtain the optimum result. Here, for three of extracted processes a
different solvent in a different ratios for both ultrasonic and soxhlet extraction (diethyl
ether, ethanol and methanol) were used respectively, but for binary solvent extraction
water and methanol were used in a different ratios, according to their polarity
extraction performed. In this study the analyte is polar, so by increasing the polarity
of solvent the extraction yield increased and gave the best result.

Determination of Extracted Amygdalin by HPLC

The mobile phase was gradient system consists of acetonitrile and deionized water
(10:90), and they gradually mix to give optimum result, the pH of the acetonitrile
was 8.41 and it was stable but the pH of water adjusted to 4.0 with 0.1M HCl. The
solvents were degassed in an ultrasonic bath prior to use.

Preparation of samples for HPLC

The extractants were prepared by dissolving the samples in an acetonitrile and then
completed to 100 ml volumetric flask.

Preparation of Standard Amygdalin

Stock solution of standard amygdalin was prepared by dissolving (1.0 g) of amyg-

dalin in methanol and then completed to 1.0 L with methanol in a volumetric flask
and stored in refrigerant before use.

Results and Discussion

Identification of Amygdalin

Extraction Methods for Amygdalin

Three extraction methods including ultrasonic extraction by different solvents (methanol

and ethanol), soxhlet extraction using the same solvents, and binary solvent extrac-
tion using (water and methanol mixture), were compared for amygdalin both sweet

119
 and bitter almond kernels. The extraction yield of amygdalin with ultrasonic extrac-
tion using methanol was the maximum, and soxhlet extraction using diethyl ether
was the minimum. Table (3) shows this result. Methanol is a good solvent for
extracting ingredients of crude herb medicine, because it is known that it can ex-
tract almost all of the ingredients, and amygdalin can be easily dissolved in it too.
Ultrasonic extraction by methanol is the best choice for the almond kernels, because
it gave the highest amount for both sweet and bitter almond kernel which are (24.45
and 563.66 µg ml−1 ) respectively.

Identification of Amygdalin with TLC

TLC test was used for identification of amygdalin from sweet and bitter almond ker-
nels, from Fig. (2) fraction (S) is the standard for amygdalin and gave a dark-brown
spot, fraction (A) gave smooth line with no spot for amygdalin extract from un-

Table 3: Amygdalin yields using different extraction methods

Yield for sweet (µg ml−1 ) Yield for bitter (µg ml−1 )
Extraction Methods
Ethanol Methanol Ethanol Methanol
Ultrasonic 22.40 24.45 525.50 563.66
Binary solvent 17.40 15.45 432.20 443.75
Soxhelt 9.98 10.44 269.83 273.21

pilled bitter almond by soxhlet extraction which indicated having many compounds
in this fraction that could not be separated well, fraction (B) gave broad dark-brown
spot which indicated that the extract containing different compounds which may be
amygdalin one of these extract from un-pilled bitter almond by ultrasonic, fraction
(C) gave one dark-brown spot which indicated that the extract containing only one
compound which was amygdalin from un-pilled bitter almond by binary solvent ex-
traction, and fraction (D) gave no spot for sweet un-pilled almond which indicated
that the extract containing no and/or less amount of amygdalin. The spots were
detected at wavelength 254 nm but were absence at 366 nm. The Rf value for the
visualization method was (0.53) for three spots of fraction S, B, and C.

Identification of Amygdalin with FTIR-Spectrum

FTIR technique was used to identify amygdaline compound in the selected samples
and compared with the standard pure amygdalin. The FTIR spectra of amygdalin in
pure form, and extracted from sweet and bitter almond kernel are shown in Figs. (3-
7). It is appeared that FTIR spectra for four extracts are relatively consistent with
the spectrum of pure amygdalin. The broad peak between 3150 and 3600 cm−1
indicates the stretching vibration of (O-H) alcoholic group in amygdalin structure.
The nitrile group (ŰC≡N) and due to the small changes in dipole moment, gives rise
to relatively low intensity band above 2359 cm−1 in all spectra. The aromatic C=C
bending appears at 1652.77 cm−1 while aromatic C-H bending at 864.11 cm−1 . A
band at 1746.26 cm−1 shows the C=0 stretching. The vibrational bands at 2854.13
cm−1 and 2925.56 cm−1 may be due to aliphatic C-H stretching vibrations and the
band at 3365.61 for C-H stretching for aromatic ring.

120
Figure 2: TLC of extract amygdalin (A-D) and standard amygdalin (S) using UV detector at 254
nm.

Figure 3: FTIR spectrum of standard pure amygdalin.

Retention time of Amygdalin from HPLC

The comparison of the retention time of the standard with those extracted from the
selected samples illustrated below, a symmetric chromatographic peak was obtained
for standard amygdaline eluted at (2.404 min). Fig. (8-10) shows the chromato-
graphic profiles for amygdalin standard and extracts of the almond sweet and bitter
kernels. It is observed that a little deviation occurred in retention times for both
types of kernels. This may be ascribed to the presence of other components in
the extracted sample matrix indicating that amygdalin is not extracted well. The
extracted amygdalin (using binary solvent extraction) from sweet and bitter almond
kernel were identified by their retention time, absorbance spectra (λmax ), parent
mass and fragmentation patterns, and where possible by comparison with pure stan-
dard amygdalin. Fig. (11-13) showed similar HPLC-DAD profiles with tR values of
(0.084, 0.079 and 0.088) min for pure amygdalin, extracted amygdalin from sweet
almond kernel and extracted amygdalin from bitter almond kernel respectively, with

121
 Figure 4: FTIR spectrum of amygdalin extracted with methanol from bitter almond kernel using
soxhlet apparatus.

Figure 5: FTIR spectrum of amygdalin extracted with methanol from sweet almond kernel using
soxhlet apparatus.

Figure 6: FTIR spectrum of amygdalin extracted with methanol from bitter almond kernel using
ultrasonic apparatus.

λmax =220 nm.

A slight change in the tR values of extracted samples compared with the tR of pure
standard may be ascribed to the type of extraction system. The fragmentation and
mass spectra data obtained with extracted amygdalin, pure standard amygdalin and
blank were shown that the peak at m/z 480 corresponds to the [M-Na]+ and m/z 475
corresponds to [M-NH4 ]+ . The peak at m/z 325 corresponds to the diglucoside ion
generated by the loss of DL-mandelonitrile [12]. The peak at m/z 172 corresponds

122
Figure 7: FTIR spectrum of amygdalin extracted with methanol from sweet almond kernel using
ultrasonic apparatus.

Figure 8: HPLC Chromatogram for standard amygdalin.

Figure 9: HPLC Chromatogram for amygdalin extract from bitter kernel by ultrasonic using
methanol solvent.

Figure 10: HPLC Chromatogram for amygdalin extract from sweet kernel by ultrasonic using
methanol solvent.

to the CH(C-2, C-6), Fig. (1) Shows the structure of amygdalin.

123
 Figure 11: HPLC/MS Chromatogram of pure standard amygdalin.

Figure 12: HPLC/MS Chromatogram of extracted amygdalin from sweet almond kernel.

Figure 13: HPLC/MS Chromatogram extracted amygdalin from bitter almond kernel.

Conclusion

The use of advanced analytical instruments such as HPLC, HPLC/MS and GC-FID
had led to identification of different phytochemical plant compounds. Analysing
these compounds by HPLC/MS has facilitated identification based on mass spectra
fragmentation pattern; however, there are limitations when it comes to identify com-
pounds with a similar molecular weight. Therefore, the need to compare absorbance
spectra and retention times (elution order) with available standards helps to verify
the identity of compounds. From this study, it is suggested that consumption of the
sweet almond kernel can be as good as the bitter almond kernel. The extraction
by ultrasonic system gave the best result (high amount) with methanol solvent for
HPLC and GC-FID, but for the HPLC-MS the extraction with binary solvent (8
ml methanol + 2 ml H2 O) were chosen for detection of amygdalin. The HPLC
method is found to be highly sensitive, reproducible, specific accurate precise and
rapid. This method can be used for the routine analysis of other extracted com-
pounds from plants. The result showed that the HPLC presented the best result than
GC-FID because the amygdalin was non volatile material.

124
Acknowledgments

We would like to express our sincere gratitude to the Kurdistan Region Government,
Ministry of Higher Education and Scientific Research and Presidency of the Uni-
versity of Salahaddin-Erbil for supporting throughout the preparation of the M.Sc.
thesis. Special thanks to KMCA (Kurdistan Medical Control Agency) and Ministry
of Natural Resource for helping and allowing us to work on their devices.

References

[1 ] SHEEL, R., NISHA, K., & KUMAR, J. (2014) ”Preliminary Phytochemical

Screening of Methanol Extract of Clerodendron infortunatum”. IOSR Journal
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[2 ] MOERTEL, C., G., FLEMING, T., R., RUBIN, J., KUOLS, L., K., SARNA,
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[3 ] YAN, J., TONG, SH., & LI, J. (2006) ”Preparative isolation and purification
of amygdalin from prunes armeniaca L. with high recovery by high-speed
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[5 ] MIAO, X., ZHAO, ZH., ZHU, H., LI, M., & ZHAO, Q. (2013) ”Compar-
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[10 ] LIANG, X., T., & FANG, W., SH. (2006) ”Medicinal Chemistry of Bioactive
Natural ProductsŤ. A John Wiley & Sons. Inc., USA.
[11 ] ARRA’ZOLA, G., SA’NCHEZ, R., DICENTA, F., & GRANE, N. (2012)
”Content of the cyanogenic glycoside amygdalin in almond seeds related to the
bitterness genotype”, Agronomia Columbiana 30 (2). p.260-265.
[12 ] ZHAO, Y. (2012) ”Amygdalin content in four stone fruit species at different
developmental stages”. Science Asia. 38. p.218-222.

125
126
 Density-Functional Calculations on Structural Properties of Sn Bulk and
Nanocrystals
Botan Jawdat Abdullah , ∗, Musafa Saeed Omar1 and Qing Jiang2
2
1
Department of Physics, College of Science, Salahaddin University, Erbil,
Kurdistan-Iraq
2
Department of Material Science and Engineering Jilin University, Changchun,
China
1
botan [email protected], 2 dr m s [email protected], 3 [email protected]
Abstract
We have performed density-functional calculations for Sn bulk and nanocrystals
using Ab initio pseudopotential method within generalize gradient approximation.
The ground-state properties are calculated for bulk lattice constant and cohesive en-
ergy, which are in good agreement with the experimental results. Our calculations
for nanocrystals analyzed. It is found that the cohesive energies diverts from that
of the bulk due to lattice constant variation. The cohesive energy decreases as the
particle sizes decreases and the effects are stronger for nanoparticles than nanowires
and nanofilms for radius less than 10nm. The results are analyzed and compared
with the available experimental data.
Keywords: Sn, Nanocrystals, Density Functional Theory, Lattice Parameterc, Co-
hesive Energy

Introduction

Low-dimension nanometer-scale materials in the form of film (NFs), wire (NWs)

and particles (NPs) are of great technological interest due to their potential appli-
cations in miniaturized electrical, optical, thermal and mechanical systems. These
nanometer-scale materials have a large surface area to volume ratio as compared to
bulk materials. Atoms at free surfaces experience a different local environment than
atoms in the bulk [1]. The fabrication of nanostructure materials and nanodevices
with desirable properties in atomic scale has become an emerging interdisciplinary
field involving solid-state physics, materials science, chemistry and biology [2].
The cohesive energy as an important physical quantity is often strongly connected
to the unique properties of nanostructures [3]. The cohesive energy of solid can
be computing by theoretical methods such as density functional theory (DFT) [4].
It has been reported that the cohesive energies measured for both Mo and W of
nanoparticles are strongly depend on their size which was carried out by Kim et al.
[5]. Currently there are several detailed theoretical models available for determining
cohesive energy of nanoparticles which are applied for metals. Jiang and his col-
leges [6] reported a model based on thermo dynamical analysis and employed well
on nanoscale size dependence of solids. DFT in the Kohn-Sham implementation
with local density approximation (LDA ) [7] or generalize gradient approximation
(GGA) [8] approximation for the exchange-correlation functional has been success-
fully applied to deduce structure, electronic, magnetic as well as other properties
of myriad of condense matter systems[9]. In the present work, the cohesive energy
based on the first-principles density functional theory calculations is investigated for
Sn bulk and nanowires (NWs) within GGA in a size range from 2 to 19 nm and the

127
 results are analyzed.

Theoretical Method

Computational Details

In the present work, the Ab-initio DFT calculations are performed by using the
ABINIT code [10] which is based on the plane-wave basis set and the total energy
self-consistently by means of first principles electronic structure is performed within
GGA using the Perdew-Burke-Ernzerh [11]. Ultrasoft Troullier-Martin [12] pseu-
dopotential was used to describe the interactions of ionic and valance electrons. In
addition, to accurately approximate integrals over the Brillion zone calculations of
the electronic states over special sets of k points, Monkhorst and pack [13] uniform
set for any crystal was used.

Lattice Constant and Cohesive Energy of Sn Bulk and Nano

The model for calculating size dependent lattice constant in bulk and nanoscale size
is reported by Omar [14] as follow:
4
a(α) = √ dmean (α) (1)
3
and (2)
4
a(r) = √ dmean (r) (3)
3
(4)

where a(α) and a(r) are the bulk and nanocrystals unit cell lattice constant and dmean
is the lattice mean bond length. The size dependent mean bond length dmean (r) is
calculated according to the following relation;

dmean (r) = h − ∆dmean (r) (5)

where h is the first surface layer inter-planer distance of the crystal. This value for
Sn is 0.3022nm, ∆dmean (r) which is the increase in mean bond length is calculated
according to:
q
−(β−1)
r
r0 −1
dmean (r) = h − ∆dmean (rC ) e (6)

where β is A related material constant and equal to 5 for Sn [14], r is the radius of
the nanowires and rC denotes a critical radius at which all atoms of the particle are
located on the bulk surface and is calculated from r0 = (3 − d)h. In this equation:
d = 0 for NPs, d = 1 for NWs and d = 2 for NFs, where r being taken as its radius
[15, 16].
The cohesive energy is an important physical quantity that accounts for the bond
strength of a solid, which equals the energy needed to divide the solid into isolated

128
atoms by breaking all the bonds. Cohesive energy is also a basic quantity for
thermodynamics, by which almost all thermodynamic properties of materials can
derive. Consider the unit cell of N atoms; the bulk cohesive energy Ecoh (α) per
unit cell is [17]:
P
Etot (α) − N Eiso
Ecoh (α) = (7)
N
where Eiso is the energy of an isolated atom. In this work, Eq. (5) is modified for
applications on nanoscale size material for GGA as follow:
P
Etot (r) − N Eiso
Ecoh = (8)
N

Results and Discussion

Different cut-off energies are tested (between 2 to 40 Hartree) and the cut-off energy
of (30 Hartree) was found to achieve better convergence within GGA. Also, accord-
ing to the Monkhorest Pack scheme the Brillion zone was sampled and (12×12×12)
mesh was used for both Sn Bulk and nanocrystals. In all cases, the energy conver-
gence was achieved with the tolerant on the difference of a total energy of less than 5
× 10-7 eV. Calculations concerning isolate energy were performed in the cubic cell
size (45A◦ ×45A◦ ×45A◦ ) with cut-off energy of (30 Hartree) to the energy conver-
gence of less than 1 × 10-11 eV. The equilibrium lattice constant value calculated
for Sn bulk structure is in good agreement compared to reported experimental data
as shown in Table (I). Calculations in light of equations from (1) to (4) give an
increase to the lattice constant, these results were used to calculate the modified
minimum total energy as well as cohesive energy for Sn NPs, NWs and NFs. The
data within GGA shows that the absolute values of minimum total energy Etot (r)
are strongly depends on the nanoscale size. Since dmean (r) for a nanowires has
large size dependence, then it has less total energy compared to bulk. On the other
hand, it has been explained that, increase in bond energy due to the size reduction
of nanocrystals is from the surface coordination imperfection formation.
Through Eq. 5, the cohesive energies Ecoh (α) were obtained by taking the energy
difference between atomic equilibrium energy in the Sn structure and the energy
of isolated atoms. The size-dependent cohesive energies for different nanocrystals
within GGA Ecoh (r) are calculated by taking the obtained minimum values of
Etot (r) due to Eq. (6) and the dependence is shown in Fig. (1) The increase of
mean bond length leads to decrease in cohesive energy with the decrease of nanosize
in this figure. Note that in the center of nanocrystals the bond length is almost the
same as it is for the bulk, which gives the bulk cohesive energy. For smaller
sizes where atoms at surface layers have larger percentage of the total number of
atoms in nanocrystals considered, the cohesive energy value is smaller. At 4 nm
sizes, Ecoh (r) reduces by 25%, 11% and 2.5% for NPs, NWs and NFs respectively.
The present calculation values on cohesive energy of Sn nanowires have the same
trend variation in comparison with the reported experimental results for Mo and W
nanoparticles [5] their absolute values of Ecoh (r) decreases with the size reductions.

129
 Table 1: Reactor and fuel element design descriptions with the HEU and LEU fuels.
Structure Properties Exp. Data This work
a(∞)/nm 6.45[18] 6.482
Ecoh (α)/eV -3.4[18] -4.2

Figure 1: The difference between the minimum total energies as a function of radius for various
nanocrystals for Sn within GGA.

Figure 2: Nanoscale size dependence cohesive energy for various nanocrystals for Sn within GGA.

Conclusions

In the present work, we studied the cohesive energy of Sn bulk and nanocrystals
from first-principles calculations based on DFT within GGA. The size dependence
of mean bond length explains well the cohesive energies of NPs, NWs and NFs.
Calculated results show that the absolute values of Ecoh (r) decreases with the de-
crease of nanoscale size and the effects are stronger for NPs than NWs and NFS for
radius less than 10nm.

Acknowledgments

The work is supported by Salahaddin-Erbil University in Kurdistan region, Iraq

under Grant No. (540-17/5/2013) as program cooperation with Jilin University in
Changchun, China.

130
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132
Geoscience Session
Assessment and Validation of Landslide Susceptibility Mapping in
Mountainous Basin Using Spatial-Based Statistical Models and Remote
Sensing Techniques
Himan Shahabi1,∗ and Bakhtyar Ali Ahmad2
1
Department of Geomorphology, Faculty of Natural Resources, University of
Kurdistan, Sanandaj, Iran
2
Department of Geoinformatics, Faculty of Geo Information and Real Estate,
Universiti Teknologi Malaysia (UTM), Malaysia
Abstract
Destructive outcomes of landslides in connection to human life along with overall
economy system of many nations around the globe are very severe. In this study,
analytical hierarchy process (AHP) model used in Expert Choice software also used
fuzzy logic with IDRISI software for prepared landslide susceptibility map for cen-
tral Zab basin in the southwest mountainsides of West-Azerbaijan province in Iran.
Ten factors as slope, aspect, elevation, lithology, NDVI, land cover, precipitation,
distance to fault, distance to drainage, and distance to road is used. Landsat ETM+
image was used for NDVI and land cover maps. Also landslide-inventory map of
the study area was identified by SPOT 5 satellite. Then fuzzy algebraic operators
were applied to the fuzzy membership values for landslide susceptibility mapping.
The fuzzy operators are validated using the receiver operating characteristic (ROC)
and the area under curve values were calculated. Based on the result of the obtained
susceptibility map by AHP model, 6.08% (31.61km2 ) of the total area show very
high landslide susceptibility. Also, among the fuzzy operators, the gamma operator
(λ = 0.975) showed the best accuracy (94.64%) while the fuzzy OR operator was
applied showed the worst accuracy (85.11%). The resultant maps would be useful
for regional spatial planning as well as for land cover planning.
Keywords: Satellite images, Landslide Susceptibility map, spatial-based statistical,
central Zab basin, Iran.

Introduction

Destructive results of landslides in human life and economy of many nations all
over the world are very severe (Nefeslioglu et al. 2008). Landslide susceptibility
is defined as a quantitative or qualitative assessment of the classification, volume
(or area), and spatial distribution of landslides which exist or may potentially occur
in an area (AGS 2007). In this respect, production of landslide susceptibility maps
at the early stage of the landslide assessments has a crucial importance for safe
and economic planning, such as urbanization activities and engineering structures,
particularly. However, a standard procedure for the production of landslide suscep-
tibility maps does not exist (Ercanoglu et al. 2004). In recent times, due to the
availability of a wide range of remote sensing data together with data from other
sources in digital form and their analysis using GIS, it has now become possible
to prepare different thematic layers corresponding to the causative factors that are
responsible for the occurrence of landslides in a region (Gupta and Joshi 1990; Van

133
 Westen 1994; Nagarajan et al. 1998; Feizizadeh et al. 2013).

The application of the analytical hierarchy process (AHP) method, developed

by Saaty (1977), has been used by many authors worldwide (Barredo et al. 2000;
Ayalew et al. 2005; Komac 2006; Yalcin 2008; Ercanoglu et al. 2008; Akgun
and Türk 2010; Shahabi et al. 2014). The advantage of fuzzy logic is that it
is straight forward to apply, and the process of weighting landslide conditioning
factors is totally controlled by the experts (Lee 2007). In addition, the fuzzy logic
method provides a variety of fuzzy combination operators for generating landslide
susceptibility index values (Ercanoglu and Gokceoglu 2002ň; Saboya Jr et al. 2006ň;
Champatiray et al. 2007; ňMuthu et al. 2008; Pradhan and Lee 2009; ňAkgun et
al. 2012; ňPourghasemi et al. 2012; Shahabi et al. 201).

The central Zab basin, in the southwest mountain sides of West-Azerbaijan

province is having problem with many landslides which occur on it. The land-
slides observed and found in the area included old and new landslides. Although
the area is not a complete agricultural zone, it has cultivated land which affected
on landslide and landslide influence on some part of it (Shahabi et al. 2012). The
main difference between this study and the references is that AHP and fuzzy re-
lations were applied for the landslide susceptibility mapping on central Zab basin,
in the southwest mountainsides of West-Azerbaijan province in Iran and landslide
causative factor databases developed using satellite images with the aid of GIS and
their comparison. However, a comparison of fuzzy approaches with another expert-
based system such as AHP is not encountered. This contribution provides originality
to this study.

Study Area Characteristics

Zab basin occupies southwestern section of West Azerbaijan and northwestern part
of Kurdistan. The area under present study covers parts of mountains and slopes
in southwestern West Azerbaijan in the central portion of Zab basin between the
latitudes of (36◦ 8’ 25”) N and (36◦ 26’ 27”) N and the longitudes of (45◦ 21’
21”) E and (45◦ 40’ 44”) E. Central Zab basin has a north-south orientation and
stretches almost 30 km in east-west direction. The study area covers some 520km2
of its total area (Figure 1). It is one of the most populated geographical basins
including a city, three towns or small cities, and over 80 villages (Khezri et al.
2013). As the study region is located in a major Zagros thrust direction, tectonic
faults are the main cause of slope failures. The regions morphology is strongly
affected by tectonic forces. The most important faults are located in western Iran,
highlighting and giving importance to this cause in the region of study. The target
research zone, tectonically, is located in Sanandaj-Sirjan region, and its north-east
region is located in Mahabad-Khoy zone (Shahabi et al. 2012).

134
 Figure 1: Geographical position of central Zab basin

Materials and Methods

Spatial database construction

In order to develop a method for the assessment of landslide susceptibility, deter-

mination of the conditioning factors for the landslides is crucial (Ercanoglu and
Gokceoglu 2002). In the first step a digital elevation model (DEM) of the study
area was generated from a triangulated irregular network (TIN) model that was
derived from digitized contours of four 1:50,000 scale topographical maps with a
contour interval of 25 m. The slope, slope aspect parameters were obtained from
the generated DEM. The critical point was the selection of appropriate pixel size
for positional accuracy and precision of susceptibility levels in the resultant map.
The positional accuracy needed for 1:50000 scale maps must be 150 m. For this
reason, a pixel size of 50 m was selected for our DEM. The DEM was then used
to describe geomorphological and geological processes in the landscape.
Fault lines were derived from 1:100000 scale geology maps of Iran Geological
organization and the aerial orthogonal distance of all pixels to fault lines calculated.
A similar process was carried out for road lines and drainage networks. In addition,
the kilometer square density of drainage networks, road and fault lines were also
used to demonstrate the importance of the features in the whole study area. Another
dataset used was land cover, which was interpreted from Landsat ETM+ image on
the 21 April 2009 that was calibrated using field observations. Due to the significant
cloud coverage, the results of the classification were edited and simplified by manual
digitization. The interpreted images were then digitally processed to further modify
the boundaries by supervision classification with ERDAS (Earth Resource Data
Analysis System) software to develop a statistical characterization of the reflectance
for each information class. Seven main land cover types were considered, namely
grassland, pasture, settled, natural forest, man-made forest, dry farm land and barren

135
 land. Based on validation from the field observations, the land cover map has
accuracy of the order of the Landsat image spatial resolution (3̃0 m).
As the factor raster maps were produced at 20 m spatial resolution, 25 m buffer
resulted in at least 4 seed cells per landslide scar. Waterways of central Zab basin
were digitized and all the needed operations for using this information in ILWIS
environment were carried out. According to the results of 30 years observations
from 3 weather stations which include Sardasht, Mirabad and Piranshahr in west
Azerbaijan province, February is the coldest month with an temperature average of
7.7 ◦ C, the hottest month is August with an temperature average of 24.1 ◦ C. In our
methodology we used long-term precipitation for a 30 years period (1980-2010).
The annual precipitation of central Zab basin is 865 mm, and the precipitations
disperse every month symmetrically. The mean annual rainfall is around 1310 mm,
most of which falls between the month of April and May. Measurements of daily
rainfall data were utilized to create average annual rainfall contour map. These were
then used for a spatial interpolation by Kriging within Arc GIS 9.3.

Landslide Susceptibility mapping models

The landslide susceptibility analyses were implemented using the AHP (analyti-
cal hierarchy process) and FL (Fuzzy Logic) methods. The maps credibility was
validated using the area under curve (AUC) and ROC method.

Analytical hierarchy process (AHP)

The AHP developed by Saaty (1980) provides a flexible and easily understood way
of analyzing complicated problems. The AHP is a multi-objective, multi-criteria
decision-making approach that allows the active participation of decision-makers in
reaching agreement, and gives managers a rational basis on which to make decisions
(Saaty and Vargas 1984). In these techniques, firstly, the effects of each parameter
to the susceptibility of landslides relative to each other were determined by dual
evaluation in determining the preferences in the effects of the parameters to the
landslide susceptibility maps (Yalcin 2008). Normally, the determination of the
values of the parameters relative to each other is a situation that is dependent on the
choices of the decision-maker. However, in this study, both the comparison of the
parameters relative to each other and the determination of the decision alternatives,
namely the effect values of the sub-criteria of the parameters (weight), were based
on the comparison of landslide inventory maps which were constructed using field
surveys, previous inventory map, and satellite image with the other data layers. The
weight of each factor was obtained from the matrix weighting factor was multiplied
by its weight class. The result of the susceptibility map is determined by factors
with high local representation such as lineaments and turned to have artifacts that
reduce its reliability. Pair-wise comparison, however, is subjective and the quality
of the results is highly dependent on the expert’s judgment. In AHP, an index of
consistency, known as the consistency ratio (CR) Eq. (1), is used to indicate the
probability that the judgments matrix were randomly generated (Saaty 1977).
CR = CI/RI (9)

136
where, RI is the average of the resulting consistency index depending on the order
of the matrix given by (Saaty 1977) and CI is the consistency index and can be
expressed as Eq. (2)

CI = (λmax − n)/(n − l) (10)

where, λmax is the largest or principal eigen value of the matrix and can be easily
calculated from the matrix and n is the order of the matrix.
For all the models, where the AHP was used, the CR was calculated. If the
CR values were greater than 0.1, the models were automatically discarded. Using
a weighted linear sum procedure (Voogd 1983) the acquired weights were used to
calculate the landslide susceptibility models (Komac 2006).

Fuzzy logic operators

The fuzzy set theory introduced by Zadeh (1965) is one of the tools used to handle
the complex problems. The fuzzy set theory employs the idea of a membership
function that expresses the degree of membership with respect to some attribute of
interest. The fuzzy logic method allows for more flexible combinations of weighted
maps, and could be readily implemented with a GIS modeling language. The idea
of using fuzzy logic in landslide hazard mapping is to consider the spatial objects
on a map as members of a set. For example, the spatial objects could be areas on an
evidence map and the set defined as areas hazardous to landslide (Pradhan and Lee
2009). Bonham-Carter (1994) discussed five operators, namely the fuzzy and, fuzzy
or, fuzzy algebraic product, fuzzy algebraic sum, and fuzzy gamma operator. This
study uses the ?ve fuzzy operators for combining the fuzzy membership functions.

Results and Discussiion

The ten factors were transformed into a vector-type spatial database using the GIS,
and landslide-related factors were extracted using the database. In this study, the
landslide susceptibility analyses were implemented using the methods of AHP and
fuzzy logic operators.

Landslide susceptibility mapping using AHP model

The result of this comparison, as a matrix was entered into the Expert choice
software. The software calculates the weight of each factor and based this weighting
factors were prioritized. The software also calculates the CR. In this study, the
CR is 0.069; the ratio indicates a reasonable level of consistency in the pair-wise
comparison that is good enough to recognize the factor weights in the landslide
susceptibility model. Precipitation with AHP weight (0.258), Lithology with AHP
weight (0.171), and slope with AHP weight (0.166) were found to be important
parameters for the study area, whereas distance to fault and distance to drainage
received a low degree of importance. The results of the pair-wise comparison
matrix, factor weights and consistency ratio of the data layers are shown in Table

137
 Table 1: Some physical and chemical properties of soil at Qlyassan

Distance
Factors to Distance Slope NDVI Lithology Aspect Distance Precipitation Elevation Land Weights
drainage to fault to road cover

Distance
to 1 0.028
drainage
Distance
1 1 0.031
to
fault
Slope 5 4 1 0.166

NDVI 2 1/2 1/5 1 0.035

Lithology 5 5 2 3 1 0.171

Aspect 3 5 1/2 4 1/2 1 0.120

Distance
2 3 1/5 3 1/2 1/3 1 0.054
to
road
Precipitation 5 6 2 5 3 3 5 1 0.258

Elevation 4 4 1/3 3 1/4 1/3 3 1/5 1 0.083

Land cover 3 3 1/5 2 1/5 1/4 2 1/5 1/2 1 0.057

Consistency ratio: 0.069

1. As a result of the AHP analyses, the landslide susceptibility map (LSM) was
produced for central Zab basin. These LSM values were then divided into four
classes based on the natural breaks range, which represent four different zones in
the landslide susceptibility. The map of landslide susceptibility has been classified
with frequency levels of 25%, 50%, and 75%.These are very high (VHS), high
(HS), moderate (MS), and low (Ayalew and Yamagishi 2005) susceptibility zones
(Fig. 2). Based on the result of the obtained susceptibility map, 19.02% (98.90km2)

Figure 2: Landslide susceptibility map generated using analytical hierarchy process for central Zab
basin

of the total area show low landslide susceptibility. Moderate, high and very high

138
Table 2: Statistics of landslide susceptibility index values in each of the ten landslide prediction
operators
NO Landslide susceptibility model LSI statistics
Min Max Mean SD
1 Fuzzy SUM 0.638 1.002 0.985 0.011
2 Fuzzy PRODUCT 0.000 0.179 0.002 0.008
3 Fuzzy GAMMA (λ=0.05) 0.000 0.198 0.006 0.010
4 Fuzzy GAMMA (λ=0.1) 0.000 0.218 0.016 0.018
5 Fuzzy GAMMA (λ=0.3) 0.000 0.305 0.020 0.025
6 Fuzzy GAMMA (λ=0.5) 0.005 0.451 0.035 0.064
7 Fuzzy GAMMA (λ=0.7) 0.110 0.631 0.076 0.089
8 Fuzzy GAMMA (λ=0.9) 0.279 0.831 0.391 0.098
9 Fuzzy GAMMA (λ=0.95) 0.353 0.911 0.426 0.085
10 Fuzzy GAMMA (λ=0.975) 0.458 0.932 0.493 0.071

susceptible zones make up 42.79% (222.50 km2), 32.11% (166.97km2) and 6.08%
(31.61km2) of the total area, respectively.

Landslide susceptibility mapping using Fuzzy Logic operators

Like the membership function, firstly the frequency ratio was calculated for each of
the landslide conditioning factors. The input factors were combined for assigning
membership functions. The statistics of the landslide susceptibility index values for
all landslide prediction operators are shown in Table 3. Using the fuzzy membership
function and the fuzzy operator, the landslide susceptibility index (LSI) values were
computed for the 10 conditioning factors are shown in Table 2. A fuzzy membership
function can provide a vehicle for developing operations with fuzzy sets in landslide
modeling. Therefore, determination of fuzzy membership values is considered to
be the most important task. The computed LSI values were mapped to allow inter-
pretation such as that illustrated, for example, in Figure 3 The values were classified
into equal areas and grouped into four classes for visual interpretation.

Validation and comparison of susceptibility maps: Success rate and predicate

rate

The landslide susceptibility analysis was performed using various fuzzy operators,
and the analysis results were validated using the Receiver Operating Characteristic
(ROC) analysis. The most ideal model shows a curve that has the largest AUC; the
AUC varying from 0.5 (random prediction, represented by the diagonal straight line)
to 1.0 (perfect prediction) (Yesilnacar and Topal 2005;ňMohammady et al. 2012).
Using the success-rate and prediction-rate methods, the 29 landslide susceptibility
maps were validated by comparing them with known landslide locations. The com-
parison results of various fuzzy operators and AHP model are shown in Figure 4.
The results show that values of AUC for the fuzzy operators vary from 86.49 to
95.64, indicating that the most of models have a reasonable good prediction capa-
bility. The fuzzy gamma operator (λ=0.975) has the highest prediction capability.
The fuzzy OR model has lowest prediction capability (Table 3). To compare the

139
 Figure 3: Landslide susceptibility maps using various fuzzy operators; a Application of fuzzy al-
gebraic OR operator; b Application of fuzzy algebraic AND operator; c Application of
fuzzy algebraic SUM operator; d Application of fuzzy algebraic PRODUCT operator; e
Application of fuzzy gamma (λ = 0.05); and f Application of fuzzy gamma (λ = 0.975).

results quantitatively, the areas under the curve were recalculated taking the total
area as 1, which means perfect prediction accuracy (Chung and Fabbri 1999). The
area under the curve is shown in Table 3. In the case of AHP model used, the area
ratio was 0.94.59 and we could say the prediction accuracy is 94.59%.

In the fuzzy operators, the area ratio for fuzzy algebraic sum was 0.9178 and
we could say that the prediction accuracy is 91.78%. In the case of applying fuzzy
algebraic product, the area ratio was 0.9459 and we could say that the prediction
accuracy is 94.59%. In the case of applying the gamma operator (λ = 0.975), the area
ratio was 0.9464 and the prediction accuracy is 94.64%. The results obtained from
Table 3 indicated that fuzzy gamma operator (λ=0.975) has the highest prediction
capability and the fuzzy OR model has lowest prediction capability (See Table 3).
Landslide susceptibility index values were visualized by means of four susceptibility
levels (low, moderate, high and very high). In this study, the equal area classification
method was used (Pourghasemi et al. 2012). Then, based on percentage of area, four
susceptibility classes were determined as very high (10%), high (20%), moderate
(30%), and low (40%). In general, the Fuzzy GAMMA (λ=0.975) model performs
better than the other models. From the prediction accuracy graphs (Fig. 6), it

140
Table 3: Statistics of landslide susceptibility index values in each of the ten landslide prediction
operators
Prediction accuracy (%)
AHP model 94.59
Fuzzy operator
Fuzzy AND 86.78
Fuzzy OR 85.11
Fuzzy algebraic sum 91.78
Fuzzy algebraic product 94.57
Fuzzy GAMMA (λ=0.05) 94.58
Fuzzy GAMMA (λ=0.1) 94.54
Fuzzy GAMMA (λ=0.3) 94.54
Fuzzy GAMMA (λ=0.5) 94.54
Fuzzy GAMMA (λ=0.7) 94.54
Fuzzy GAMMA (λ=0.9) 94.58
Fuzzy GAMMA (λ=0.95) 94.59
Fuzzy GAMMA (λ=0.975) 94.64

is quite evident that the susceptibility map with fuzzy gamma operator shows the
similar prediction accuracy of 95%.

Figure 4: Success rate curve evaluation of landslide susceptibility maps produced by fuzzy operators

Conclusion

In the present study, landslide susceptibility in a landslide prone area located in the
southwest mountainsides of West-Azerbaijan province in Iran was assessed using
fuzzy operators and AHP model by considering ten landslide conditioning factors.
The known landslides were then overlaid the landslide susceptibility maps to examine
the degree of coincidence of susceptibility with land sliding. After validation, among

141
 the 12 operators, the case of applying the gamma operator (λ =0.975), showed
the best accuracy (94.64%), whereas the fuzzy algebraic OR (85.11%) and fuzzy
AND (86.78%) operators showed the worst accuracy. In the case of applying the
gamma operator with different λ value, the prediction accuracy had a similar upward
trend value, between 94.58 and 94.64%. Generally, the validation results showed
satisfactory agreement between the susceptibility map and the existing data from
landslide locations. The effect of choosing different values of gamma (between 0
and 1) is not large. Because the landslide susceptibility maps using different values
of gamma are very similar and the prediction accuracy after validation is also very
similar. In the AHP model, the precipitation (weight= 0.258) and lithology (weight=
0.171) parameters are positively associated with the occurrence of landslides. The
verification process also showed that the predictive capability of the method is
acceptable. Thus, the spatial distributions of landslide susceptibility in the resulting
maps by the AHP model and fuzzy operators were in?uenced by the precipitation
factor. To confirm the practicality of the results, the three susceptibility maps were
compared with 29 active landslide zones by using ROC method.
The validation results by ROC method shows that the area under the curve for
AHP model and fuzzy gamma operator (?=0.975) are 0.9459 (94.59%) and 0.9464
(94.64%) with prediction accuracy 94%, and 94% respectively. The verification
process also showed that the predictive capability of the method is acceptable.
So, the map extracted from the fuzzy gamma operator (λ=0.975) combined with
the remote sensing and GIS data, yields a reasonable accuracy for the landslide
prediction and higher accuracy as compared to maps extracted from the other fuzzy
operators and AHP model. Areas within the high and very high susceptibility
categories should require further study by engineering geologists before development
to determine the extent of possibly unstable conditions. The information provided by
this landslide susceptibility map could be the basis for decision making, as planners
and engineers can reduce losses caused by existing and future landslides by means
of prevention, mitigation and avoidance.

Acknowledgments

We are thankful to the Department of Geomorphology, Faculty of Natural Resources,

University of Kurdistan, Iran for providing the facilities for this investigation.

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146
Mathematics Session
αγ -Connectedness and Some Properties of α(γ,β) -Continuous Functions
ALIAS B. KHALAF1 AND HARIWAN Z. IBRAHIM2
1
Department of Mathematics, Faculty of Science, University of Duhok,
Kurdistan-Region, Iraq
2
Department of Mathematics, Faculty of Science, University of Zakho, Kurdistan-
Region, Iraq
[email protected] and hariwan [email protected]
Abstract
The aim of this paper first is to introduce the notion of αγ -separated sets and study
their properties in topological spaces, then we introduce the notions of αγ -connected
and αγ -disconnected spaces. We discuss the characterizations and properties of
αγ -connected spaces and then properties under α(γ,β) -continuous functions. The
αγ -components in a space X and αγ -locally connected spaces are also introduced.
Keywords: operation, αγ -open, αγ -closed, αγ -closure, αγ -connected and α(γ,β) -
continuous functions.

Introduction

The notion of connectedness is useful and fundamental not only in general topology
but also in other advanced branches of mathematics. Many researchers have inves-
tigated the basic properties of connectedness. The productivity and fruitfulness of
connectedness motivated mathematicians to generalize this notion. In the course of
these attempts many stronger and weaker forms of connectedness have been intro-
duced and investigated. In 1965, Njastad introduced the notion of an α-open set in
a topological space, and proved that the collection of all α-open sets in (X, τ ) is a
topology on X, finer that τ . Having two related topologies on the same underlying
set, it is quite natural to ask whether they share some topological properties. Several
authors expanded the well-known notion of Connectedness [3] to some topological
spaces. In [7], [8] and [9] they studied different types of connectedness. In 2013,
Ibrahim have introduced the concept αγ -open set in topological spaces and investi-
gated its properties. In the present paper, we shall introduce and study another type
of connectedness by means of operations defined on the α-topology of the space.
Throughout this paper (X, τ ) and (Y, σ) are topological spaces with no separation
axioms assumed unless otherwise stated. For a subset A of a space X, Cl(A) and
Int(A) denote the closure and interior of A, respectively.

Preliminaries

In this section, we recall some of the basic definitions and some important results.
Definition 1 (5). A subset A of a topological space (X, τ ) is called α-open if
A ⊆ Int(Cl(Int(A))).

147
 The complement of an α-open set is called α-closed. The intersection of all
α-closed sets containing A is called the α-closure of A and is denoted by αCl(A).
By αO(X), we denote the family of all α-open sets of X.

Definition 2 (4). An operation γ : αO(X) → P (X) is a mapping satisfying the

condition, V ⊆ V γ for each V ∈ αO(X). We call the mapping γ an operation on
αO(X). The operation id : αO(X) → P (X) is defined by V id = V for any set
V ∈ αO(X) this operation is called the identity operation on αO(X).

By γ and β we always mean operations defined on the family of all α-open sets
of the space. We recall the following definitions and results from [4].
Definition 3. 1. A subset A of X is called an αγ -open set if for each point
x ∈ A, there exists an α-open set U of X containing x such that U γ ⊆ A.
The complement of an αγ -open set is said to be αγ -closed. We denote the set
of all αγ -open sets of (X, τ ) by αO(X)γ .
2. The αγ -closure of a subset A of X with an operation γ on αO(X) is denoted by
αγ Cl(A) and is defined to be the intersection of all αγ -closed sets containing
A.
3. A mapping f : (X, τ ) → (Y, σ) is said to be α(γ,β) -continuous if for each x
of X and each αβ -open set V containing f (x), there exists an αγ -open set U
such that x ∈ U and f (U ) ⊆ V .
4. A mapping f : (X, τ ) → (Y, σ) is said to be α(γ,β) -closed if for any αγ -closed
set A of (X, τ ), f (A) is αβ -closed in (Y, σ).
Theorem 4. A point x ∈ X is in the αγ -Closure of the subset A if and only if every
αγ -open set V of X containing x has a non-empty intersection with A.
Recalling from [1], that the αγ -interior of A is the union of all αγ -open sets
contained in A and it is denoted by αγ Int(A). The αγ -boundary of a set A is
defined as αγ Cl(A) \ αγ Int(A) and is denoted by αγ Bd(A).
Theorem 5. [1] If A is a subset of X, then αClγ (X \ A) = X \ αIntγ (A).
Definition 6 (6). The subsets A and B of a topological space (X, τ ) are called
α-τ -separated sets if (αCl(A) ∩ B) ∪ (A ∩ αCl(B)) = φ.
Definition 7 (2). A mapping f : (X, τ ) → (Y, σ) is said to be α(γ,β) -open if for any
αγ -open set A of (X, τ ), f (A) is αβ -open in (Y, σ).
Corollary 8 (2). A mapping f : (X, τ ) → (Y, σ) is α(γ,β) -closed if and only if
αβ Cl(f (A)) ⊆ f (αγ Cl(A)) for every subset A of (X, τ ).

αγ -Connected Spaces

Definition 9. Two subsets A and B of a topological space (X, τ ) are called αγ -

separated if (αγ Cl(A) ∩ B) ∪ (A ∩ αγ Cl(B)) = φ.
Remark 10. Each two αγ -separated sets are always disjoint, since A ∩ B ⊆ A ∩
αγ Cl(B) = φ. The converse may not be true in general.

148
Example 11. Let X = {a, b, c} and τ = {φ, X, {a}}. Define an operation γ on
αO(X) by 
γ A if c ∈ A
A =
X if c ∈
/ A.
Then, {a} and {c} are disjoint subsets of X but not αγ -separated.
Remark 12. From the fact that αCl(A) ⊆ αγ Cl(A), for every subset A of X. Then
every αγ -separated set is α-τ -separated. But the converse may not be true as shown
in the following example.
Example 13. Let X = {a, b, c, d} and τ = {φ, X, {a}, {b}, {a, b}}. Define an
operation γ on αO(X) by

γ A if d ∈ A
A =
X if d ∈ / A.

Then, the subsets {c} and {d} are α-τ -separated but not αγ -separated.
Theorem 14. If A and B are any two nonempty subsets in a space X, then the
following statements are true:
1. If A and B are αγ -separated, A1 ⊆ A and B1 ⊆ B, then A1 and B1 are also
αγ -separated.
2. If A ∩ B = φ such that each of A and B are both αγ -closed (αγ -open), then
A and B are αγ -separated.
3. If each of A and B are both αγ -closed (αγ -open) and if H = A ∩ (X \ B) and
G = B ∩ (X \ A), then H and G are αγ -separated.
Proof
1. Since A1 ⊆ A, then αγ Cl(A1 ) ⊆ αγ Cl(A). Then, B ∩ αγ Cl(A) = φ implies
B1 ∩ αγ Cl(A) = φ and B1 ∩ αγ Cl(A1 ) = φ. Similarly A1 ∩ αγ Cl(B1 ) = φ.
Hence, A1 and B1 are αγ -separated.
2. Since A = αγ Cl(A) and B = αγ Cl(B) and A∩B = φ, then αγ Cl(A)∩B = φ
and αγ Cl(B) ∩ A = φ. Hence, A and B are αγ -separated. If A and B are
αγ -open, then their complements are αγ -closed. Hence, αγ Cl(A) ⊆ X \ B
and αγ Cl(B) ⊆ X \ A. Therefore, A and B are αγ -separated.
3. If A and B are αγ -open, then X \ A and X \ B are αγ -closed. Since H ⊆
X \ B, αγ Cl(H) ⊆ αγ Cl(X \ B) = X \ B and so αγ Cl(H) ∩ B = φ. Thus
G ∩ αγ Cl(H) = φ. Similarly, H ∩ αγ Cl(G) = φ. Hence H and G are αγ -
separated. If A and B are αγ -closed, then αγ Cl(H) ⊆ A and αγ Cl(G) ⊆ B.
Thus, H and G are αγ -separated.
Theorem 15. The sets A and B of a space X are αγ -separated if and only if there
exist U and V in αO(X)γ such that A ⊆ U , B ⊆ V and A ∩ V = φ, B ∩ U = φ.
Proof Let A and B be αγ -separated sets. Set V = X \ αγ Cl(A) and U =
X \ αγ Cl(B). Then U, V ∈ αO(X)γ such that A ⊆ U , B ⊆ V and A ∩ V = φ,
B ∩ U = φ. On the other hand, let U, V ∈ αO(X)γ such that A ⊆ U , B ⊆ V
and A ∩ V = φ, B ∩ U = φ. Since X \ V and X \ U are αγ -closed, then
αγ Cl(A) ⊆ X \V ⊆ X \B and αγ Cl(B) ⊆ X \U ⊆ X \A. Thus αγ Cl(A)∩B = φ
and αγ Cl(B) ∩ A = φ.

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 Theorem 16. In any topological space (X, τ ), the following statements are equiv-
alent:
1. φ and X are the only both αγ -open and αγ -closed sets in X.
2. X is not the union of two disjoint nonempty αγ -open sets.
3. X is not the union of two disjoint nonempty αγ -closed sets.
4. X is not the union of two nonempty αγ -separated sets.
Proof (1) ⇒ (2): Suppose (2) is false and that X = A ∪ B, where A, B are
disjoint, nonempty and αγ -open. Since X \ A = B is αγ -open and nonempty, we
have that A is a nonempty proper αγ -open and αγ -closed set in X, which shows
that (1) is false.
(2) ⇔ (3): This is clear.
(3) ⇒ (4): If (4) is false, then X = A ∪ B, where A, B are nonempty and
αγ -separated. Since αγ Cl(B) ∩ A = φ, we conclude that αγ Cl(B) ⊆ B, so B is
αγ -closed. Similarly, A must be αγ -closed. Therefore, (3) is false.
(4) ⇒ (1): Suppose (1) is false and that A is a nonempty proper αγ -open and
αγ -closed subset of X. Then, B = X \ A is nonempty, αγ -open and αγ -closed, so
A and B are αγ -separated and X = A ∪ B, so (4) is false.
Definition 17. A subset C of a space X is said to be αγ -disconnected if there are
nonempty αγ -separated subsets A and B of X such that C = A ∪ B, otherwise C
is called αγ -connected. If C is αγ -disconnected, such a pair of sets A, B will be
called an αγ -disconnection of C.
Example 18. Let X = {a, b, c} and τ = {φ, X, {a}, {b}, {a, b}, {b, c}}. Define an
operation γ on αO(X) by

γ A if c ∈ A
A =
Cl(A) if c ∈ / A.
Then, X is αγ -disconnected because there exist a pair {a}, {b, c} subsets of X
such that {a} ∪ {b, c} = X, and (αγ Cl({a}) ∩ {b, c}) ∪ ({a} ∩ αγ Cl({b, c})) =
({a} ∩ {b, c}) ∪ ({a} ∩ {b, c}) = φ.
Example 19. Let X = {a, b, c} and τ = {φ, X, {a}, {c}, {a, c}}. Define an opera-
tion γ on αO(X) by 
γ A if b ∈ A
A =
X if b ∈ / A.
Then, X is αγ -connected, since there does not exist a pair A, B of nonempty αγ -
separated subsets of X such that X = A ∪ B.
Remark 20.
• Every indiscrete space is αγ -connected.
• Every discrete space contains more than one element is αid -disconnected.
• A space X is αγ -connected if any of the conditions (1) − (4) in Theorem 16
hold.
• According to Definition 17 and Remark ??, a space X is αγ -disconnected if
we can write X = A ∪ B, where one of the following statements is true:

150
 1. A and B are disjoint, nonempty and αγ -open.
2. A and B are disjoint, nonempty and αγ -closed.
3. A and B are nonempty and αγ -separated.
Theorem 21. A space X is αγ -disconnected if and only if there exists a nonempty
proper subset A of X which is both αγ -open and αγ -closed in X.
Proof Follows from Remark 20.
Theorem 22. A space X is αγ -connected if and only if every nonempty proper
subset of X has a nonempty αγ -boundary.
Proof Suppose that a nonempty proper subset A of an αγ -connected space X
has empty αγ -boundary. Since αγ Cl(A) = αγ Int(A) ∪ αγ Bd(A). Thus, A is both
αγ -closed and αγ -open. By Theorem 21, X is αγ -disconnected. This contradiction,
hence proves that A has a nonempty αγ -boundary.
Conversely, suppose X is αγ -disconnected. Then by Theorem 21, X has a
nonempty proper subset A which is both αγ -closed and αγ -open. Then, αγ Cl(A) =
A, αγ Cl(X \ A) = X \ A and αγ Cl(A) ∩ αγ Cl(X \ A) = φ. So A has empty
αγ -boundary, this is a contradiction. Hence, X is αγ -connected.
Lemma 23. Suppose M, N are αγ -separated subsets of X. If C ⊆ M ∪ N and C
is αγ -connected, then C ⊆ M or C ⊆ N .
Proof Since C ∩ M ⊆ M and C ∩ N ⊆ N , then C ∩ M and C ∩ N are αγ -
separated and C = C ∩ (M ∪ N ) = (C ∩ M ) ∪ (C ∩ N ). But C is αγ -connected
so (C ∩ M ) and (C ∩ N ) can not form an αγ -disconnection of C. Therefore, either
C ∩ M = φ, so C ⊆ N or C ∩ N = φ, so C ⊆ M .
Theorem 24. Suppose C and Ci (i ∈ I) are αγ -connected subsets of X and that
for each i, Ci and C are not αγ -separated. Then, S = C ∪ Ci is αγ -connected.
Proof Suppose that S = M ∪ N , where M and N are αγ -separated. By Lemma
23, either C ⊆ M or C ⊆ N . Without loss of generality, assume C ⊆ M . By the
same reasoning we conclude that for each i, either Ci ⊆ M or Ci ⊆ N . But if some
Ci ⊆ N , then C and Ci would be αγ -separated. Hence every Ci ⊆ M . Therefore,
N = φ and the pair M, N is not an αγ -disconnection of S.
Corollary 25. Suppose that for each i ∈ I, Ci is an αγ -connected subset of X and
that for all i 6= j, Ci ∩ Cj 6= φ. Then, ∪{Ci : i ∈ I} is αγ -connected.
Proof If I = φ, then ∪{Ci : i ∈ I} = φ is αγ -connected. If I 6= φ, pick i0 ∈ I
and let Ci0 be the central set C in Theorem 24. For all i ∈ I, Ci ∩ Ci0 6= φ, so Ci
and Ci0 are not αγ -separated. By Theorem 24, ∪{Ci : i ∈ I} is αγ -connected.
Corollary 26. Suppose that for all x, y ∈ X, there exists an αγ -connected set
Cxy ⊆ X with x, y ∈ Cxy . Then, X is αγ -connected.
Proof Certainly X = φ is αγ -connected. If X 6= φ, choose a ∈ X. By hypothesis
there is, for each y ∈ X, an αγ -connected set Cay containing both a and y. By
Corollary 16, X = ∪{Cay : y ∈ X} is αγ -connected.
Corollary 27. Suppose C is an αγ -connected subset of X and A ⊆ X. If C ⊆ A ⊆
αγ Cl(C), then A is αγ -connected.
For each a ∈ A, {a} and C are not αγ -separated. By Theorem 24,
Proof S
A = C ∪ {{a} : a ∈ A} is αγ -connected.

151
 Remark 28. In particular, the αγ -closure of an αγ -connected set is αγ -connected.
Corollary 29. A function f : X → Y is α(γ,β) -continuous if and only if f −1 (V ) is
αγ -open in X, for every αβ -open set V in Y .
Proof Let V be an αβ -open set in Y . Then, Y \ V is an αβ -closed set in
Y . Then by [[4], Theorem 4.2 (2)], f −1 (Y \ V ) is αγ -closed set in X. Since
f −1 (Y \ V ) = X \ f −1 (V ) and hence f −1 (V ) is αγ -open set in X.
Conversely, Let x ∈ X and V be an αβ -open subset of Y containing f (x).
Then, x ∈ f −1 (V ). By hypothesis f −1 (V ) is αγ -open in X contains x and clearly
f (f −1 (V )) ⊆ V . Hence, f is α(γ,β) -continuous.
Theorem 30. If f : (X, τ ) → (Y, σ) is an onto α(γ,β) -continuous function and X is
αγ -connected, then Y is αβ -connected.
Proof Suppose that Y is αβ -disconnected and A, B is an αβ -disconnection of Y .
By Remark 20, A and B are both αβ -open sets. Since f is α(γ,β) -continuous, so by
Corollary 29, f −1 (A) and f −1 (B) are both nonempty αγ -open sets in X. Now,
f −1 (A) ∪ f −1 (B) = f −1 (A ∪ B) = f −1 (Y ) = X, and
f −1 (A) ∩ f −1 (B) = f −1 (A ∩ B) = f −1 (φ) = φ.
Then by Remark 20, f −1 (A), f −1 (B) is a pair of αγ -disconnection of X. This
contradiction shows that Y is αβ -connected.
Theorem 31. Let f : (X, τ ) → (Y, σ) be an injective function, then following are
equivalent:
1. f is α(γ,β) -continuous.
2. f −1 (V ) ⊆ αγ Int(f −1 (V )) for every αβ -open set V of Y .
3. f (αγ Cl(A)) ⊆ αβ Cl(f (A)) for every subset A of X.
4. αγ Cl(f −1 (B)) ⊆ f −1 (αβ Cl(B)) for every subset B of Y .
5. f −1 (αβ Int(B)) ⊆ αγ Int(f −1 (B)) for every subset B of Y .
Proof (1) ⇒ (2): Let V be any αβ -open set of Y and x ∈ f −1 (V ). Then, f (x) ∈
V . Since f is α(γ,β) -continuous, there exists an αγ -open set U of X containing x
such that f (U ) ⊆ V and hence U ⊆ f −1 (V ), this implies that x ∈ αγ Int(f −1 (V )).
Thus, it follows that f −1 (V ) ⊆ αγ Int(f −1 (V )).
(2) ⇒ (3): Let A be any subset of X and f (x) ∈ / αβ Cl(f (A)). Then, by
Theorem 4, there exists an αβ -open set V of Y containing f (x) such that V ∩f (A) =
φ and hence f −1 (V ) ∩ A = φ. Also f (x) ∈ V implies x ∈ f −1 (V ). Then by (2)
we obtain that x ∈ αγ Int(f −1 (V )). Hence, there exists an αγ -open set U of X
containing x such that U ⊆ f −1 (V ). Then U ∩ A = φ and so x ∈ / αγ Cl(A). This
implies f (x) ∈
/ f (αγ Cl(A)). Thus, f (αγ Cl(A)) ⊆ αβ Cl(f (A)).
(3) ⇒ (4): Let B be any subset of Y . Since f (f −1 (B)) ⊆ B, so, we
have αβ Cl(f (f −1 (B))) ⊆ αβ Cl(B). Also, f −1 (B) ⊆ X, then by (3), we have
f (αγ Cl(f −1 (B))) ⊆ αβ Cl(f (f −1 (B))) ⊆ αβ Cl(B). Thus, αγ Cl(f −1 (B)) ⊆
f −1 (αβ Cl(B)).
(4) ⇒ (5): Let B be any subset of Y and x ∈ f −1 (αβ Int(B)). Then,
/ X \ f −1 (αβ Int(B)) = f −1 (αβ Cl(Y \ B)). By (4), x ∈
by Theorem 5, x ∈ /

152
αγ Cl(f −1 (Y \ B)) = X \ αγ Int(f −1 (B)) and hence x ∈ αγ Int(f −1 (B)). Thus,
f −1 (αβ Int(B)) ⊆ αγ Int(f −1 (B)).
(5) ⇒ (1): Let x ∈ X and V be any αβ -open set of Y containing f (x).
Since V ∩ Y \ V = φ, we have f (x) ∈ / αβ Cl(Y \ V ) = Y \ αβ Int(V ) and
hence f (x) ∈ αβ Int(V ) which implies x ∈ f −1 (αβ Int(V )). By (5), we obtain that
x ∈ αγ Int(f −1 (V )). This means that there exists an αγ -open set U of X containing
x such that U ⊆ f −1 (V ) and so f (U ) ⊆ V . This shows that f is α(γ,β) -continuous.
Corollary 32. Let f : X → Y be an α(γ,β) -continuous and injective function. If K
is αγ -connected in X, then f (K) is αβ -connected in Y.
Proof
Suppose that f (K) is αβ -disconnected in Y . There exist two αβ -separated sets
P and Q of Y such that f (K) = P ∪Q. Set A = K ∩f −1 (P ) and B = K ∩f −1 (Q).
Since f (K) ∩ P 6= φ, then K ∩ f −1 (P ) 6= φ and so A 6= φ. Similarly B 6= φ.
Now, A ∪ B = (K ∩ f −1 (P )) ∪ (K ∩ f −1 (Q)) = K ∩ (f −1 (P ) ∪ f −1 (Q)) =
K ∩ f −1 (P ∪ Q) = K ∩ f −1 (f (K)) = K. Since f is α(γ,β) -continuous, then
by Theorem 31 (4), αγ Cl(f −1 (Q)) ⊆ f −1 (αβ Cl(Q)) and B ⊆ f −1 (Q), then
αγ Cl(B) ⊆ f −1 (αβ Cl(Q)). Since P ∩ αβ Cl(Q) = φ, then A ∩ αγ Cl(B) ⊆
A ∩ f −1 (αβ Cl(Q)) ⊆ f −1 (P ) ∩ f −1 (αβ Cl(Q)) = φ and then A ∩ αγ Cl(B) = φ.
Similarly, B ∩ αγ Cl(A) = φ. Thus, A and B are αγ -separated. Therefore, K is
αγ -disconnected, this is contradiction. Hence, f (K) is αβ -connected.
Corollary 33. For a bijective α(γ,β) -closed function f : X → Y , if C is αβ -
connected in Y , then f −1 (C) is αγ -connected in X.
Proof Suppose that f −1 (C) is αγ -disconnected in X. There exist two αγ -
separated sets M and N of X such that f −1 (C) = M ∪ N . Set K = C ∩ f (M )
and L = C ∩ f (N ). Since C = f (M ) ∪ f (N ), then C ∩ f (M ) 6= φ and so K 6= φ.
Similarly L 6= φ. Now, K ∪ L = (C ∩ f (M )) ∪ (C ∩ f (N )) = C ∩ (f (M ) ∪ f (N )) =
C ∩f (M ∪N ) = C ∩f (f −1 (C)) = C. Since f is α(γ,β) -closed, then by Corollary 8,
αβ Cl(f (N )) ⊆ f (αγ Cl(N )) and L ⊆ f (N ), then αβ Cl(L) ⊆ f (αγ Cl(N )). Since
M ∩αγ Cl(N ) = φ, then K∩αβ Cl(L) ⊆ K∩f (αγ Cl(N )) ⊆ f (M )∩f (αγ Cl(N )) =
φ and then K ∩ αβ Cl(L) = φ. Similarly, L ∩ αβ Cl(K) = φ. Thus, K and L are αβ -
separated. Therefore, C is αβ -disconnected, this is contradiction. Hence, f −1 (C)
is αγ -connected.
Theorem 34. A space X is αγ -disconnected if and only if there exists an α(γ,id) -
continuous function f from X onto the discrete space {0, 1}.
Proof Suppose that X is αγ -disconnected. Then, there exist two nonempty dis-
joint αγ -open subsets G1 and G2 of X such that X = G1 ∪ G2 . Define a function
f : X → {0, 1} as follows.

0 if x ∈ G1
f (x) =
1 if x ∈ G2 .

Now, the only αid -open sets in {0, 1} are φ, {0}, {1} and {0, 1}. So, f −1 (φ) =
φ, f −1 ({0}) = G1 , f −1 ({1}) = G2 and f −1 ({0, 1}) = X, which are αγ -open sets in
X. Thus, by Corollary 29, f is α(γ,id) -continuous surjection from X to the discrete
space {0, 1}.

153
 Conversely, let the hypothesis holds and if possible suppose that X is αγ -
connected. Therefore, by Theorem 30, {0, 1} is αid -connected, which is contra-
diction by Remark 20. So, X must be αγ -disconnected.
Theorem 35. A space X is αγ -connected if and only if every α(γ,id) -continuous
function from X to the discrete space {0, 1} is constant.
Proof Let X be αγ -connected and f : X → {0, 1} be any α(γ,id) -continuous
function. Let y ∈ f (X) ⊆ {0, 1}, then {y} ⊆ {0, 1} and since {0, 1} is discrete,
so {y} is both αid -open and αid -closed in {0, 1}. Since f is α(γ,id) -continuous.
Therefore, by Corollary 29 and [[4], Theorem 4.2 (2)], f −1 ({y}) is both αγ -open
and αγ -closed in X. Now, since y ∈ f (X). So, there exists x ∈ X such that
y = f (x). Thus, f (x) ∈ {y} and x ∈ f −1 ({y}). Hence, we obtain f −1 ({y}) 6= φ.
If f −1 ({y}) 6= X, then f −1 ({y}) is a nonempty subset of X which is both αγ -open
and αγ -closed. By Theorem 21, which implies that X is αγ -disconnected, this is a
contradiction, so f −1 ({y}) = X. Thus, f (X) = {y}, it means that f (x) = y, for
each x ∈ X, so f is constant.
Conversely, let the hypothesis be holds and if possible suppose that X is αγ -
disconnected. Therefore, by Theorem 21, X has a nonempty proper subset A of
X which is both αγ -open and αγ -closed. So, X \ A is also a nonempty proper
subset of X which is both αγ -open and αγ -closed. Now, consider the characteristic
function ΨA of A defined as

0 if x ∈ A
ΨA (x) =
1 if x ∈ X \ A.
Now, Ψ−1 −1 −1 −1
A (φ) = φ, ΨA ({0}) = A, ΨA ({1}) = X \ A and ΨA ({0, 1}) = X, which
are all αγ -open sets in X. So, ΨA is α(γ,id) -continuous function from X to the
discrete space {0, 1}. By hypothesis, ΨA must be constant, this is contradiction,
since ΨA is not constant function. So, X is αγ -connected.
Definition 36. A set C is called maximal αγ -connected set if it is αγ -connected and
if C ⊆ D ⊆ X where D is αγ -connected, then C = D. A maximal αγ -connected
subset C of a space X is called an αγ -component of X. If X is itself αγ -connected,
then X is the only αγ -component of X.
Theorem 37. For each x ∈ X, there is exactly one αγ -component of X containing
x.
S
Proof For any x ∈ X, let Cx = {A : x ∈ A ⊆ X and A is αγ -connected}.
Then, {x} ∈ Cx , since Cx is a union of αγ -connected sets each containing x, Cx is
αγ -connected by Corollary 16. If Cx ⊆ D and D is αγ -connected, then D was one
of the sets A in the collection whose union defines Cx , so D ⊆ Cx and therefore
Cx = D. Therefore, Cx is an αγ -component of X that contains x.
Corollary 38. A space X is the union of its αγ -components.
Proof Follows from Theorem 37.
Corollary 39. Two αγ -components either are disjoint or coincide.
Proof Let Cx and Cy be αγ -components and Cx 6= Cy . If p ∈ Cx ∩ Cy , then
by Corollary 16, Cx ∪ Cy would be an αγ -connected set strictly larger than Cx .
Therefore, Cx ∩ Cy = φ.

154
Theorem 40. Each αγ -connected subset of X is contained in exactly one αγ -
component of X.
Proof Let A be an αγ -connected subset of X which is not in exactly one αγ -
component of X. Suppose that C1 and C2 are αγ -components of X such that
A ⊆ C1 and A ⊆ C2 . Since C1 ∩ C2 6= φ and by Corollary 16, C1 ∪ C2 is another
αγ -connected set which contains C1 as well as C2 , a contradiction to the fact that
C1 and C2 are αγ -components. This proves that A is contained in exactly one
αγ -component of X.
Theorem 41. A nonempty αγ -connected subset of X which is both αγ -open and
αγ -closed is αγ -component.
Proof Suppose that A is αγ -connected subset of X which is both αγ -open and
αγ -closed. By Theorem 40, A is contained in exactly one αγ -component C of X. If
A is a proper subset of C, then C = (C∩A)∪(C∩(X \A)) and (C∩A), (C∩(X \A))
is an αγ -disconnection of C, which is a contradiction. Thus, A = C.
Theorem 42. Every αγ -component of X is αγ -closed.
Proof Suppose that C is an αγ -component of X. Then, by Remark 28, αγ Cl(C)
is αγ -connected containing αγ -component C of X. This implies that C = αγ Cl(C)
and hence C is αγ -closed.
Definition 43. A space X is called αγ -locally connected if for each point x ∈ X
and for each αγ -open set U containing x, there is an αγ -connected αγ -open set V
such that x ∈ V ⊆ U .
Theorem 44. Let a function f : X → Y be α(γ,β) -continuous, α(γ,β) -open and
bijective. If X is αγ -locally connected, then Y is αβ -locally connected.
Proof Let y ∈ Y and choose x ∈ X such that f (x) = y. Let U be an αβ -open
set of Y containing y. Since f is α(γ,β) -continuous, then f −1 (U ) is αγ -open by
Corollary 29 and also since X is αγ -locally connected, then there exists an αγ -
connected αγ -open set V containing x such that x ∈ V ⊆ f −1 (U ). This gives that
f (x) ∈ f (V ) ⊆ f f −1 (U ) = U or y ∈ f (V ) ⊆ U . Since f is α(γ,β) -open, so f (V )
is αβ -open. Moreover f (V ) is αβ -connected by Corollary 32. This proves that Y
is αβ -locally connected.
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156
 Maximal λc -Open Sets
Halgwrd M. Darwesh1 , Sarhad F. Namiq2 and Wria K. Kadir1
1
Department of Mathematics, School of Science, Faculty of Science and Science
Education, University of Sulaimani, Sulaimani, Kurdistan-Iraq
2
Department of Mathematics, School of Science, Faculty of Education, University
of Garmyan, Kalar-Sulaimani Sulaimani, Kurdistan-Iraq.
[email protected], [email protected] and
[email protected]
Abstract
The goal of the present paper is to introduce and define a new class of maximal
open sets which is called maximal λc -open sets. Some of it’s basic properties,
characterization and relationships are studies and obtained.
Keywords: s-operation λ, λc -open sets, maximal open sets, maximal λc - open sets.

Preliminaries and Introduction

The notion of semi-open set were introduced by Levine [4] in 1963. Nakaoka and
Oda [1] in 2003 studied the idea of maximal open sets, then after analyzing this
concept, they introduced the concept of minimal closed-sets [2]. However, Namiq
[3] defined an s-operation λ and used it to define and study the concept of λc -open
sets. In the current paper, we introduce and discuss the notion of maximal λc -open
sets. If X is a topological space and A ⊆ X then cl(A) and int(A) are used to
denote the closure and interior of A respectively. A subset A of a topological space
X is said to be semi-open [4] if A ⊆ cl(int(A)). The complement of a semi-open
set is called semi-closed. The family of all semi-open of a topological space X will
denote by SO(X). The function λ : SO(X) → P (X), where P (X) is the power set
of X is called an s-operation [3] if A ⊆ λA for each A ∈ SO(X), and its assumed
λφ = φ always.
In what follows, by a space X, we mean a topological space (X, τ ) with an
s-operation λ. The following definitions are recalled from [3].
Definition 1. A semi-open subset A of a space X is said to be λc -open if for each
x ∈ X, there exist a semi-open set U and a closed set F both containing x such
that λU and F ⊆ A.
The complement of a λc -open set is called λc -closed and the family of all λc -open
subsets of a space X is denoted by λc O(X).
Definition 2. For a subset A of a space X,
1. λc cl(A) = ∪{F ; F is λc − closed s.t A ⊆ F } is the λc -closure of A.
2. λc int(A) = ∩{U ; U is λc − open s.t. U ⊆ A} is the λc -interior of A.
Theorem 3. The uniion of any family of λc -open subsets of X is λc -open.
Theorem 4. For a subset A and a point x of a space X, we have:
1. x ∈ λc cl(A) iff G ∩ A 6= φ, for each λc -open subset G of X containing x.
2. x ∈ λc int(A) iff there exists a λc -open set G such that x ∈ G ⊆ A.
Theorem 5. For subsets A and B of a space X, we have:

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 1. If A ⊆ B then
• λc cl(A) ⊆ λc cl(B)
• λc int(A) ⊆ λc int(B)
2. A is λc -closed iff A = λc cl(A)
3. A is λc -open iff A = λc int(A)
4. • X \ λc int(A) = λc cl(X \ A)
• X \ λc cl(A) = λc int(X \ A)
Definition 6. A non-empty proper open subset A of a space X is said to be maximal
open [1] if any open subset of X contains A is A or X. The complement of a maximal
open set is minimal closed [2].

Maximal λc -Open Sets

In this section, we introduce and study the notion of maximal λc -open sets.
Definition 7. A non-empty proper λc -open subset A of a space X is said to be
maximal λc -open if any λc -open set contains A is either A or X.
It is clear that, every maximal λc -open set is λc -open, and hence its semi-open.
But not conversely, it can be seen in the following example.
Example 8. Let X = {a, b, c} be a set equipped with the discrete ( topology τd . Then
A A ∈ {φ, {a}, {a, b}}
SO(X) = P (X). Now, if we define λ : SO(X) → P (X) by λA = ,
X o.w
then the set {a} is λc -open but it is not maximal λc -open.
By the following example, we show that the notion of maximal open and maximal
λc -open are independent:
Example 9. Let (X, τd ) be( the space of Example 8. Define an s-operation λ :
A A ∈ {φ, {a}}
SO(X) → P (X) by λA = , then the set A = {a} is maximal
X o.w
λc -open but not maximal open. However, the set B = {a, b} is maximal open but it
is not maximal λc -open.
Proposition 10. Let X be any space, Then:
1. If A is a maximal λc -open subset of X and B is λc -open subset of X then
either B ⊆ A or A ∪ B = X.
2. If A and B are maximal λc -open subsets of X, then either A ∪ B = X or
A = B.
Proof
1. Since A and B are λc -open subsets of X, then by Theorem 3, A ∪ B is a
λc -open subset of X contains A. So by Definition 7 either A ∪ B = A or
A ∪ B = X. That is either B ⊆ A or A ∪ B = X.
2. It follows from part 1.

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Proposition 11. Let A , B and C be three maximal λc -open subsets of a space X
such that A 6= B. If A ∪ B ⊆ C then either C = A or C = B.
Proof Since A and C are maximal λc -open subsets of X, then by the part 2 of
the Proposition 10 either A ∪ C = X or C = A. If A ∪ C = X, then

B = B ∩ X = B ∩ (A ∪ C) = (B ∩ A) ∪ (B ∩ C) ⊆ C.

Since B is a maximal λc -open set, then by the part 2 of the Proposition 10 either
C = B or C = X. Since C 6= X, so C = B. Hence we have either C = A or
C = B.
Proposition 12. Let A be a maximal λc -open subset of a space X and let x be a
point of A. Then

A = ∪{B : B is a λc -open set contains x and A ∪ B 6= X}.

Proof Let C = {B : B is a λc -open set contains X and A ∪ B =6= X}. Since

A itself is a λc -open set contains x and A ∪ A = A 6= X, so A ⊆ ∪C. Conversely,
for each B ∈ C, B is a λc -open set in which x ∈ B and A ∪ B 6= X. By part 1 of
Proposition 10, we have B ⊆ A, so ∪C ⊆ A. Therefore A = ∪C.
Proposition 13. If A is a maximal λc -open subset of a space X and x ∈ X \ A,
then A ∪ G = X for each λc -open set G contains x.
Proof Let G be any λc -open subset of X such that x ∈ G. Since A is maximal
λc -open, then by part 1 of the Proposition 10 either A ∪ G = X or G ⊆ A. But
since G 6⊆ A, so A ∪ G = X.
Corollary 14. If A is a maximal λc -open subset of a space X and x ∈ X \ A then:
1. X \ A ⊆ G, for each λc -open subset G of X contains x.
2. X \ G ⊆ A, for each λc -open subset G of X contains x.
3. H ⊆ A, for each λc -closed subset H of X not containing x.
Proof It follows from Proposition 13.
Proposition 15. Let A be a non-empty proper λc -open subset of a space X such that
each point of X \ A is contained in a finite λc -closed set. Then there is a maximal
λc -open subset B of X such that A ⊆ B and X \ B is finite.
Proof Since A is a non-empty proper subset of X, then there is an element
x ∈ X \ A, so by our hypothesis, there is a finite set λc -closed set F such that
x ∈ F . By Theorem 3, A ∪ X \ F is a λc -open subset of X contains A. If A is a
maximal λc -open, then by the part 1 of the Proposition 10 either A ∪ X \ F = X or
X \ F ⊆ A. Since x ∈ X \ F ∩ A, then X \ F ⊆ A so X \ A ⊆ F which is finite.
Thus the set B = A is the required set. If A ∪ X \ F is maximal λc -open, then the
set B = A ∪ X \ F is the required set. Otherwise, there exists a λc -open set A0 ,
such that A ∪ X \ F & A0 & X If A0 is maximal λc -open set, then the set B = A0
is the required set (since X \ B = X \ A0 ⊆ F which is finite). If A0 is not maximal
λc -open, then there exist a λc -open set A1 such that A ∪ X \ F & A0 & A1 & X.
Following this process gives a sequence of λc -open sets Am , m ∈ N such that

A ∪ X \ F & A0 & A1 & A2 & ... & Am & ... & X

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 Since X \ A ∩ F ⊆ F which is finite, so this sequence can not be infinite. That
is, there is n ∈ N such that Am = An for all m ≥ n. So the set B = An is the
required set.
Corollary 16. Let A be a non-empty proper cofinite λc -open subset of a space X.
Then there exists a cofinite maximal λc -open set B contains A.
Proof It follows from Proposition 15.
The existence of the maximal λc -open set in Proposition 15 and Corollary 16 is
not unique, as the following example shows:
a space of Example 8 and let λ : SO(X) → P (X) be an
Example 17. Let (X, τd ) be (
A A ∈ {φ, {a}, {a, b}, {a, c}}
s-operation given by λA = . Then B = {a, b}
X o.w
and C = {a, c} are maximal λc -open sets in X containing the cofinite set A = {a}.
Proposition 18. Either λc cl(A) = A or λc cl(A) = X, for any maximal λc -open
subset A of a space X.
Proof If λc cl(A) = A, then there is nothing to prove. Otherwise, to show
λc cl(A) = X. On the contrary, we suppose that λc cl(A) 6= X, this means that, there
is a point x ∈ X\λc cl(A). Then by Theorem 4, there exists a λc -open set G contains
x such that G ∩ A = φ. Since A is a maximal λc -open set, then by the part 1 of
the Proposition 10 either A ∪ G = X or G ⊆ A. But since G * A, so A ∪ G = X.
This implies that A = X \ G. By Theorem 5, λc cl(A) = λc cl(X \ G) = X \ G = A
which is impossible. Hence λc cl(A) = X.
Proposition 19. Let A be a maximal λc -open subset of a space X and let B be any
non-empty subset of X \ A. Then λc cl(B) = X \ A.
Proof Since X \ A is λc -closed and B ⊆ X \ A, then by Theorem 5, we have
λc cl(B) ⊆ X \ A. To show X \ A ⊆ λc cl(B), let x ∈ / λc cl(B). Then by Theorem 4,
there exists a λc -open set G such that x ∈ G and G ∩ B = φ. Since A is a maximal
λc -open set, then by part 1 of the Proposition 10, either A ∪ G = X or G ⊆ A. If
A ∪ G = X, then B = B ∩ X = B ∩ (A ∪ G) = (B ∩ A) ∪ (B ∩ G) = φ which is a
contradictions to the hypothesis that B 6= φ. So, we have G ⊆ A, this implies that
x ∈ A. Therefore, x ∈ / X \ A. Thus λc cl(B) = X \ A.
Proposition 20. If A is a maximal λc -open subset of a space X, then λc int(C) = A,
for each proper superset C of A.
Proof Since A ⊆ C, then by Theorem 5, A = λc int(A) ⊆ λc int(C). Since C is
a proper subset of X and λc int(C) ⊆ C, then λc int(C) is a proper λc -open subset
of X contains A. Since A is maximal λc -open, so λc int(C) = A.
Proposition 21. Let A and Aγ be maximal λc -open subsets of a space X for each
γ ∈ Γ such that ∩ Aγ ⊆ A. Then there exists γ ∈ Γ such that A = Aγ .
γ∈Γ

Proof If A 6= Aγ for each γ ∈ Γ, then by part 2 of the Proposition 10 A∪Aγ = X

for each γ ∈ Γ. That is, ∩ (A ∩ Aγ ) = X, so that A ∩ ( ∩ Aγ ) = X. This implies
γ∈Γ γ∈Γ
that A = X which is impossible. Hence A = Aγ for some γ ∈ Γ.
Corollary 22. Let {Aγ }γ∈Γ be a family of maximal λc -open subsets of a space X.
If A is also a maximal λc -open subset of X such that A 6= Aγ for each γ ∈ Γ, then:

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 1. X \ ∩ Aγ ⊆ A
γ∈Γ

2. ∩ Aγ ∩ A = φ
γ∈Γ

Proof Both followed from Pproposition 21

Proposition 23. Let {Aγ }γ∈Γ be a family of maximal λc -open subsets of a space
X with Γ has more than one element. If Aγ 6= Aσ for each γ 6= σ in Γ, then
Aσ = ( ∩ Aγ ) ∪ (X \ ∩ Aγ ) for all σ ∈ Γ.
γ∈Γ γ∈Γ\{σ}

Proof Let σ ∈ Γ. Then by Corollary 15, (X \ ∩ Aγ ) ⊆ Aσ = A , so

γ∈Γ\{σ}

A = A ∩ X = A ∩ [( ∩ Aγ ) ∪ (X \ ∩ Aγ )]
γ∈Γ\{σ} γ∈Γ\{σ}

= (A ∩ ∩ Aγ ) ∪ (X \ ∩ Aγ ) = ∩ Aγ ∪ (X \ ∩ Aγ ).
γ∈Γ\{σ} γ∈Γ\{σ} γ∈Γ γ∈Γ\{σ}

Corollary 24. Let {Aγ }γ∈Γ be a family of maximal λc -open subsets of a space
X with Γ has more than two elements. If Aγ 6= Aσ for each γ 6= σ in Γ, then
Aγ ∩ Aσ 6= φ for each γ, σ ∈ Γ.
Proof It follows from Proposition 23.
Finally we obtain some results on relationship between maximal open sets and
λc -open sets.
Proposition 25. A non-empty λc -open subset A of a space X is maximal open if
and only if there exists a maximal open set G such that G ⊆ A ⊆ cl(G).
Proof If A is maximal open, then the set G = A is the required maximal open
set. Conversely, Suppose that G is the maximal open set such that G ⊆ A ⊆ cl(G),
then by Definition 6 the set X \ G is minimal closed. If (X \ G) ∩ A 6= φ, then there
is a point x ∈ X \ G and x ∈ A. So by Definition 7, there exists a closed set F
such that x ∈ F ⊆ A. By [[2] , part 1 of Lemma 2.1] X \ G ⊆ F ⊆ A. Therefore
X = G ∪ X \ G ⊆ A ∪ A = A which is impossible (since A is proper subset of X).
Hence (X \ G) ∩ A = φ, so A ⊆ G. Thus A = G which is maximal open.
Corollary 26. A non-empty proper λc -open subset A of a space X is maximal open
if and only if there exists a maximal open set G such that G ⊆ A ⊆ scl(G), where
scl(G) is the semi-closure of G in X.
Proof It follows from Proposition 20 and the fact that scl(G) ⊆ cl(G).
Corollary 27. If A is a non-open maximal λc -open subset of a space X, then int(A)
is not maximal open.
Proof It follows from Proposition 20 and the fact that every λc -open set is semi-
open.
Corollary 28. Let A be a maximal λc -open subset of a space X. Then int(A) is
maximal open if and only if int(A) = A.
Proof It follows from Proposition 25 and Corollary 27.
Corollary 29. Let G be a maximal open subset of a space X. If A is any λc -open
subset of X such that G ⊆ A, then either A = G or A = X.

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 Proof Since G ⊆ A, then G ⊆ int(A). Since G is maximal open, then either
int(A) = G or int(A) = X. If int(A) = X, then A = X 6= G. If int(A) = G, then
G ⊆ A ⊆ cl(G), so by Proposition 25 A = G 6= X.
Example 9 shows that the converse of Corollary 29 is not true ingeneral.
Corollary 30. Let A be a λc -open subset of X which contains a maximal open set,
then either cl(A) = A or cl(A) = X.
Proof From Corollary 29 either A = G or A = X, and hence either cl(A) =
cl(G) = G = A or cl(A) = cl(G) = X.
Definition 31. A non-empty proper λc -closed subset A of a space X is called
minimal λc -closed if any λc -closed set which is contained in A is either φ or A.
It is easy to show the following:
Proposition 32. A non-empty proper subset A of X is minimal λc -closed if and only
if X \ A is maximal λc -open.
Remark 33. In view of Proposition 32 and Theorem 5 and the fact that (A ⊆ B
if and only if X \ B ⊆ X \ A). It is easy to establish the dual results of minimal
λc -closed sets as we stated them on maximal λc -open sets.

References
[1 ] F. Nakaoka and N. Oda, Some Properties of Maximal Open-Sets, Interna-
tional Journal of Mathematics.Math.Sci.21, 2003, pp. 1331-1340.
[2 ] F. Nakaoka and N. Oda, Minimal Closed Sets and Maximal Closed Sets,
International Journal of Mathematics.Math.Sci.21, 2006, Article ID 1-8.
[3 ] S.F. Namiq, New Types of Continuity and Separation Axioms Based on
Operation in Topological Spaces,M.Sc. Thesis, College of Science, University
of Sulaimani, 2011.
[4 ] N. Levine, Semi-Open Sets and Semi-Continuity in Topological Spaces,
AMER. Math.Soc.70, 1963, pp. 957–964.

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 New R0 and R1 in Bitopological spaces
Hardi N. Aziz
Department of Mathematics, School of Science Education, Faculty of Science and
Science Education, University of Sulaimani, Sulaimani, Kurdistan-Iraq
[email protected]
Abstract
In this paper, introduce two new classes of bitopological spaces called (i, j) −
SC − R0 and (i, j) − SC − R1 in terms of the concept (i, j) − SC −open set and
investigate some of their fundamental properties.
Keywords: Bitopological spaces, (i, j) − SC −open set, (i, j) − SC − R0 , (i, j) −
SC − R 1 .

Introduction

The separation axioms R0 and R1 where introduced and studied by N. A. Shan

in [9] and C. T. yang in [2]. In 1963 they were rediscovered by A. S. Davis [1],
Then, their some properties are obtained by some authors [6, 7]. In [5], Kelly was
first introduced the concepts of bitopological spaces, where X is a non-empty set
and τ1 , τ2 are topologies on X. In [8] Levine introduced the concepts of semi-
open sets in topological spaces. In this paper, we investigate further properties
of (i, j) − SC − R0 and (i, j) − SC − R1 , throughout this paper, (X, τ1 , τ2 ) is a
bitopological spaces. However, to avoid cumbersome expressions, let i and j refer
to the first and second topologies on X, respectively.

Preliminaries

The following results can be found in [3]

Definition 1. A subset A of a bitopological space (X, τ1 , τ2 ) is said to be (i, j)−SC −
open, if A is j− semi- open and for all x ∈ A, there exist an i− closed set F such
that x ∈ F ⊆ A. A subset B of X is called (i, j) − SC − closed if and only if
X − B is (i, j) − SC − open. The family of (i, j) − SC − open (resp., (i, j) − SC −
closed) subset of X is denoted by (i, j) − SC O(X) (resp., (i, j) − SC C(X)).
Definition 2. Let (X, τ1 , τ2 ) be a bitopological space and x ∈ X. A subset N of X
is said to be (i, j) − SC − neighborhood of x, if there exists an (i, j) − SC − open
set U in X such that x ∈ U ⊆ N .
Definition 3. If A is a subset of a bitopological space (X, τ1 , τ2 ), then the (i, j) −
SC − closure of A is the intersection of all (i, j)−SC − closed sets which containing
A, and denoted by (i, j) − SC Cl(A).
The following definition can be found in [4]
Definition 4. A bitopological spaces (X, τ1 , τ2 ) is said to be (i, j) − SC − T0 , if for
every distinct points x and y of X, there exists an (i, j) − SC − open set containing
one of them but not the other.
Definition 5. A bitopological spaces (X, τ1 , τ2 ) is said to be (i, j) − SC − T1 , if
for every distinct points x and y of X, there exist two (i, j) − SC − open sets, one
containing x but not y and the other containing y but not x.

163
 Definition 6. A bitopological spaces (X, τ1 , τ2 ) is said to be (i, j) − SC − T2 , if for
every distinct points x and y of X, there exist two disjoint (i, j) − SC − open sets,
one containing x and the other containing y.
Definition 7 (8). A subset A of a topological space (X, τ ) is said to be semi- open,
if A ⊆ Cl(Int(A)), denoted by SO(X). The complement of semi- open sets is
called semi- closed set, denoted by SC(X).
Theorem 8. [4] For any bitopological space (X, τ1 , τ2 ), the following statements
are equivalent:
1- (X, τ1 , τ2 ) is (i, j) − SC − T0 .
2- (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}), for all x 6= y in X.
3- Either x ∈
/ (i, j) − SC Cl({y}) or y ∈
/ (i, j) − SC Cl({x}), for all x 6= y in X.
4- A relation R defined on X × X such that for all x 6= y, xRy if and only if
x ∈ (i, j) − SC Cl({y}) is anti- symmetric.

(i, j) − SC − R0 and (i, j) − SC − R1

Definition 9. A bitopological space (X, τ1 , τ2 ) is called (i, j)−SC −R0 , if for every
(i, j) − SC − open set A, (i, j) − SC Cl({x}) ⊆ A, for all x ∈ A.
Definition 10. A bitopological space (X, τ1 , τ2 ) is called (i, j)−SC −R1 , if for each
x, y ∈ X, (i, j) − SC Cl(x) 6= (i, j) − SC Cl({y}), there exist disjoint (i, j) − SC −
open sets U and V such that (i, j) − SC Cl({x}) ⊆ U and (i, j) − SC Cl(y) ⊆ V .
Definition 11. Let (X, τ1 , τ2 ) be a bitopological space and A ⊆ X. Then the
(i, j) − SC − Kernal of A denoted by (i, j) − SC − ker(A) and defined by the
intersection of all (i, j) − SC − open sets which containing A.
Lemma 12. Let (X, τ1 , τ2 ) be a bitopological space and x, y ∈ X. Then y ∈
(i, j) − SC − ker({x}) if and only if x ∈ (i, j) − SC Cl({y}). Proof Suppose that
y∈ / (i, j) − SC − ker({x}), then there exists an (i, j) − SC − open set V containing
x such that y ∈ / V . Therefore, x ∈ / (i, j) − SC Cl({y}). The proof of converse can
be done similarly.
Example 13. Let X = {a, b, c}, τ1 = {∅, {a}, {b, c}, X}, τ2 = {∅, {b}, X}. Then
the space (X, τ1 , τ2 ) is (i, j) − SC − R0 space.
Theorem 14. For any bitopological space (X, τ1 , τ2 ) the following properties are
equivalent:
1- (X, τ1 , τ2 ) is (i, j) − SC − R0 .
2- For any (i, j) − SC − closed set F and x ∈ X, there exists an (i, j) − SC −
open set U such that F ⊆ U and x ∈ / U.
3- For any (i, j) − SC − closed set F with x ∈
/ F implies that F ∩ (i, j) −
SC Cl({x}) = ∅.
4- For any two distinct points x, y ∈ X, either (i, j) − SC Cl({x}) = (i, j) −
SC Cl({y}) or (i, j) − SC Cl({x}) ∩ (i, j) − SC Cl({y}) = ∅.

164
 Proof
1 =⇒ 2: Let F ∈ (i, j) − SC C(X) and x ∈ / F . Then by (1), (i, j) − SC Cl({x}) ⊆
X − F . Set U = X − [(i, j) − SC Cl({x})], then U ∈ (i, j) − SC O(X), F ⊆ U
and x ∈ / U.
2 =⇒ 3: Let F ∈ (i, j) − SC C(X) and x ∈ / F . Then there exists an (i.j) − SC −
open set U such that x ∈ U , since F ⊆ U, so by (2), U ∩ (i, j) − SC Cl({x}) = ∅,
implies that F ∩ (i, j) − SC Cl({x}) = ∅.
3 =⇒ 4: For any distinct points x, y ∈ X, suppose that (i, j)−SC Cl({x}) 6= (i, j)−
SC Cl({y}), then there exist z ∈ (i, j)−SC Cl({x}) such that z ∈ / (i, j)−SC Cl({y})
and there exists V ∈ (i, j) − SC O(X) such that y ∈ / V and z ∈ V , hence x ∈ V .
Therefore x ∈ / (i, j) − SC Cl({y}). By (3), we obtain (i, j) − SC Cl({x}) ∩ (i, j) −
SC Cl({y}) = ∅.
4 =⇒ 1: Let V ∈ (i, j) − SC O(X) and x ∈ V , for each y ∈ / V . Then x 6= y and
x∈ / (i, j) − SC Cl({y}), this shows that (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}).
By (4), (i, j) − SC Cl({x}) ∩ (i, j) − SC Cl({y}) = ∅, for each y ∈ X − V . Hence
(i, j) − SC Cl({x}) ∩ [∪(i, j) − SC Cl({y})] = ∅, where y ∈ X − V . On the other
hand, since V ∈ (i, j) − SC O(X) and y ∈ X − V, so (i, j) − SC Cl({y}) ⊆ X − V .
Hence, X − V = ∪(i, j) − SC Cl({y}), where y ∈ X − V . Therefore, obtain that
(i, j) − SC Cl({x}) ∩ (X − V ) = ∅, so (i, j) − SC Cl({x}) ⊆ V . Hence, (X, τ1 , τ2 )
is (i, j) − SC − R0 space.
Proposition 15. For any two distinct points x, y in X, the following statements are
equivalent:
1- (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}).
2- (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}).
Proof
1 =⇒ 2: Suppose that (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}). Then there
exists a point z ∈ X such that z ∈ (i, j) − SC − ker({x}) and z ∈ / (i, j) − SC −
ker({y}). Since z ∈ (i, j) − SC − ker({x}), so {x} ∩ (i, j) − SC Cl({z}) 6= ∅,
this implies that x ∈ (i, j) − SC Cl({z}), and by z ∈ / (i, j) − SC − ker({y}),
we have {y} ∩ (i, j) − SC Cl({z}) = ∅. Since x ∈ (i, j) − SC Cl({z}), so (i, j) −
SC Cl({x}) ⊆ (i, j)−SC Cl({z}) and hence {y}∩(i, j)−SC Cl({x}) = ∅. Therefore,
(i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}).
2 =⇒ 1: Suppose that (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}). Then there exists
a point z ∈ X such that z ∈ (i, j) − SC Cl({x}) and z ∈ / (i, j) − SC Cl({y})
and there exists an (i, j) − SC − open set containing z and hence x but not y, so
y∈ / (i, j)−SC −ker({x}). Therefore, (i, j)−SC −ker({x}) 6= (i, j)−SC −ker({y}).
Theorem 16. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − R0 if and only if for
any points x, y ∈ X, (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}) implies that
(i, j) − SC − ker({x}) ∩ (i, j) − SC − ker({y}) = ∅.
Proof
Necessity: Suppose that (X, τ1 , τ2 ) is (i, j) − SC − R0 . Thus by Proposition 15
for any points x, y in X, if (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}), then
(i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}). To prove that (i, j) − SC − ker({x}) ∩
(i, j) − SC − ker({y}) = ∅. Assume that z ∈ (i, j) − SC − ker({x}) ∩ (i, j) −
SC − ker({y}), since z ∈ (i, j) − SC − ker({x}), so by Lemma 12 it follows

165
 that x ∈ (i, j) − SC Cl({z}). Since x ∈ (i, j) − SC Cl({x}), so by Theorem 14
(i, j) − SC Cl({x}) = (i, j) − SC Cl({z}), similarly we have (i, j) − SC Cl({y}) =
(i, j) − SC Cl({z}) = (i, j) − SC Cl({x}), which is contradiction. Therefore, (i, j) −
SC − ker({x}) ∩ (i, j) − SC − ker({y}) = ∅.
Sufficiency: Let (X, τ1 , τ2 ) be a bitopological space such that for any points x and
y in X, (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}) implies that (i, j) − SC −
ker({x}) ∩ (i, j) − SC − ker({y}) = ∅. If (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}),
hence by Proposition 15 (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}). Therefore
(i, j)−SC −ker({x})∩(i, j)−SC −ker({y}) = ∅, implies that (i, j)−SC Cl({x})∩
(i, j) − SC Cl({y}) = ∅, because, if z ∈ (i, j) − SC Cl({x}) implies that x ∈
(i, j) − SC − ker({z}), so (i, j) − SC − ker({x}) ∩ (i, j) − SC − ker({z}) 6= ∅,
by hypothesis we have (i, j) − SC − ker({x}) = (i, j) − SC − ker({z}), then
z ∈ (i, j) − SC Cl({x}) ∩ (i, j) − SC Cl({y}), implies that(i, j) − SC − ker({x}) =
(i, j) − SC − ker({z}) = (i, j) − SC − ker({y}), which is contradiction. Therefore,
(i, j) − SC Cl({x}) ∩ (i, j) − SC Cl({y}) = ∅, by Theorem 14 (X, τ1 , τ2 ) is (i, j) −
SC − R0 space.
Theorem 17. For any bitopological space (X, τ1 , τ2 ) the following statements are
equivalent:
1- (X, τ1 , τ2 ) is (i, j) − SC − R0 space.
2- For any nonempty set A and G ∈ (i, j) − SC O(X) such that A ∩ G 6= ∅, there
exists F ∈ (i, j) − SC C(X) such that A ∩ F 6= ∅ and F ⊆ G.
3- For any G ∈ (i, j) − SC O(X), G = ∪{F ∈ (i, j) − SC C(X) : F ⊆ G}.
4- For any F ∈ (i, j) − SC C(X), F = ∩{G ∈ (i, j) − SC O(X) : F ⊆ G}.
5- For any x ∈ X, (i, j) − SC Cl({x}) ⊆ (i, j) − SC − ker({x}).
Proof
1 =⇒ 2: Let A be a nonempty subset of X and G ∈ (i, j) − SC O(X) such that
A ∩ G 6= ∅. Then there exists a point x ∈ X such that x ∈ A ∩ G, since x ∈ G ∈
(i, j) − SC O(X), so by (1), (i, j) − SC Cl({x}) ⊆ G. Set F = (i, j) − SC Cl({x}),
then F ∈ (i, j) − SC C(X), F ⊆ G and A ∩ F 6= ∅.
2 =⇒ 3: Let G ∈ (i, j) − SC O(X). Then ∪{F ∈ (i, j) − SC C(X) : F ⊆ G} ⊆ G
and let x be any point of G. By (2), there exists F ∈ (i, j) − SC C(X) such that
x ∈ F and F ⊆ G. Therefore, x ∈ F ⊆ ∪{F ∈ (i, j) − SC C(X) : F ⊆ G}. Hence
G = ∪{F ∈ (i, j) − SC C(X) : F ⊆ G}.
3 =⇒ 4: Clear.
4 =⇒ 5: Let x be any point of X and y ∈ / (i, j) − SC − ker({x}). So there exists an
(i, j)−SC − open set V such that x ∈ V and y ∈ / V , hence (i, j)−SC Cl({y})∩V =
∅. By (4),[∩{G ∈ (i, j) − SC O(X) : (i, j) − SC Cl({y}) ⊆ G}] ∩ V = ∅, where
G ∈ (i, j) − SC O(X) such that x ∈ / G and (i, j) − SC Cl({y}) ⊆ G. Therefore,
(i, j) − SC Cl({x}) ∩ G = ∅ and hence y ∈ / (i, j) − SC Cl({x}). Consequently we
obtain (i, j) − SC Cl({x}) ⊆ (i, j) − SC − ker({x}).
5 =⇒ 1: Let G ∈ (i, j) − SC O(X), with x ∈ G and let y ∈ (i, j) − SC − ker({x}).
Then x ∈ (i, j)−SC Cl({y}) and y ∈ G, this implies that (i, j)−SC −ker({x}) ⊆ G.
Therefore, x ∈ (i, j) − SC Cl({x}) ⊆ (i, j) − SC − ker({x}) ⊆ G. Hence (X, τ1 , τ2 )
is (i, j) − SC − R0 space.

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Proposition 18. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − R0 space if and
only if (i, j) − SC Cl({x}) = (i, j) − SC − ker({x}), for all x ∈ X.
Proof
Necessity: Suppose that (X, τ1 , τ2 ) is (i, j) − SC − R0 . By Theorem 17, (i, j) −
SC Cl({x}) ⊆ (i, j)−SC −ker({x}), for all x ∈ X. Let y ∈ (i, j)−SC −ker({x}),
then x ∈ (i, j) − SC Cl({y}) and by Theorem 14, (i, j) − SC Cl({x}) = (i, j) −
SC Cl({y}). Therefore, y ∈ (i, j) − SC Cl({x}) and hence, (i, j) − SC − ker({x}) ⊆
(i, j) − SC Cl({x}). Hence, (i, j) − SC Cl({x}) = (i, j) − SC − ker({x}), for all
x ∈ X.
Sufficiency: Follows From Theorem 17.
Proposition 19. For any bitopological space (X, τ1 , τ2 ) the following statements
are equivalent:
1) (X, τ1 , τ2 ) is (i, j) − SC − R0 .
2) x ∈ (i, j) − SC Cl({y}) if and only if y ∈ (i, j) − SC Cl({x}), for any points
x, y in X.
Proof
1 =⇒ 2: Assume that (X, τ1 , τ2 ) is (i, j) − SC − R0 . Let x ∈ (i, j) − SC Cl({y})
and A be any (i, j) − SC − open set containing y. By (1), (i, j) − SC Cl({y}) ⊆ A,
hence x ∈ A. Therefore, every (i, j) − SC − open set containing y contains x, so
y ∈ (i, j) − SC Cl({x}).
2 =⇒ 1: Let G be any (i, j)−SC − open set containing x, if y ∈ / G, then x ∈/ (i, j)−
SC Cl({y}). By (2), y ∈ / (i, j) − SC Cl({x}), this implies that (i, j) − SC Cl({x}) ⊆
G. Therefore, (X, τ1 , τ2 ) is (i, j) − SC − R0 .
Corollary 20. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − T0 if and only if
(i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}), for all x and y in X. Proof
Follows from Theorem 8 and Proposition 15.
Proposition 21. Every (i, j) − SC − T1 is (i, j) − SC − R0 .
Proof Obvious. The following result shows that the converse of Proposition
21 is not true in general:
Example 22. Consider the space in Example 13, we see that X is (i, j) − SC − R0
but not (i, j) − SC − T1 .
Theorem 23. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − T1 if and only if it
is an (i, j) − SC − T0 and (i, j) − SC − R0 space.
Proof
Necessity: Let X be (i, j) − SC − T1 . Then by Proposition 21 X is (i, j) − SC − R0
and since every (i, j) − SC − T1 is (i, j) − SC − T0 [4], which completes the proof.
Sufficiency: Let X is both (i, j)−SC −T0 and (i, j)−SC −R0 space. To show X is
(i, j)−SC −T1 . Let x, y ∈ X be any two distinct points. Since X is (i, j)−SC −T0 ,
so there exists an (i, j) − SC − open set G such that x ∈ G and y ∈ / G or there
exists an (i, j) − SC − open set H such that y ∈ H and x ∈ / H. Suppose that x ∈ G
and y ∈ / G, since X is (i, j) − SC − R0 , then (i, j) − SC Cl({x}) ⊆ G, as y ∈ /G
implies that y ∈ / (i, j) − SC Cl({x}). Hence, y ∈ H = X − [(i, j) − SC Cl({x})], it
clear that x ∈/ H. So there exist (i, j) − SC − open sets G and H containing x and
y, respectively such that y ∈ / G and x ∈/ H. Hence, X is (i, j) − SC − T1 .

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 Theorem 24. Let (X, τ1 , τ2 ) be a bitopological space, A ⊆ X and x ∈ X. Then
(i, j) − SC − ker(A) = {x ∈ X : (i, j) − SC Cl({x}) ∩ A 6= ∅}.
Proof Let x ∈ (i, j) − SC − ker(A) and (i, j) − SC Cl({x}) ∩ A = ∅. Then
x∈/ X −[(i, j)−SC Cl({x})] which is an (i, j)−SC − open set containing A. This is
impossible, because x ∈ (i, j)−SC −ker(A). Therefore, (i, j)−SC Cl({x})∩A 6= ∅.
On the other hand, let (i, j) − SC Cl({x}) ∩ A 6= ∅ and x ∈/ (i, j) − SC − ker(A),
then there exist an (i, j) − SC − open set U containing A such that x ∈ A, let
y ∈ (i, j) − SC Cl({x}) ∩ A, so U is (i, j) − SC − neighborhood of y which x ∈/ U,
which is contradiction. Hence, x ∈ (i, j) − SC − ker(A).
Theorem 25. For any bitopological space (X, τ1 , τ2 ) the following statements are
equivalent:
1- (X, τ1 , τ2 ) is (i, j) − SC − R0 .
2- If F is (i, j) − SC − closed, then F = (i, j) − SC − ker(F ).
3- If F is (i, j) − SC − closed and x ∈ F , then (i, j) − SC − ker({x}) ⊆ F .
4- If x ∈ X, then (i, j) − SC − ker({x}) ⊆ (i, j) − SC Cl({x}).
Proof
1 =⇒ 2: Let F be an (i, j) − SC − closed set and x ∈ / F . Thus X − F is an
(i, j)−SC − open set contains x. Since X is (i, j)−SC −R0 , so (i, j)−SC Cl({x}) ⊆
X − F , thus (i, j) − SC Cl({x}) ∩ F = ∅, by Theorem 24, x ∈ / (i, j) − SC − ker(F ),
hence (i, j) − SC − ker(F ) ⊆ F . Therefore, F = (i, j) − SC − ker(F ).
2 =⇒ 3: In general, if A ⊆ B, then (i, j) − SC − ker(A) ⊆ (i, j) − SC − ker(B).
Therefore, by (2), (i, j) − SC − ker({x}) ⊆ (i, j) − SC − ker(F ) = F .
3 =⇒ 4: Since x ∈ (i, j) − SC Cl({x}) and (i, j) − SC Cl({x}) is (i, j) − SC −
closed set, so by (3), (i, j) − SC − ker({x}) ⊆ (i, j) − SC Cl({x}).
4 =⇒ 1: Let x ∈ (i, j)−SC Cl({y}). Then by Lemma 12, y ∈ (i, j)−SC −ker({x}),
since x ∈ (i, j) − SC Cl({x}) and (i, j) − SC Cl({x}) is (i, j) − SC − closed set,
so by (4), y ∈ (i, j) − SC − ker({x}) ⊆ (i, j) − SC Cl({x}). Therefore, x ∈
(i, j) − SC Cl({y}) implies that y ∈ (i, j) − SC Cl({x}), where x, y ∈ X. Hence by
Proposition 19, (X, τ1 , τ2 ) is (i, j) − SC − R0 .
Definition 26. A bitopological space (X, τ1 , τ2 ) is said to be (i, j) − SC − R1 , if
for each x, y ∈ X, (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}), there exist disjoint
(i, j) − SC − open sets U and V such that (i, j) − SC Cl({x}) ⊆ U and (i, j) −
SC Cl({y}) ⊆ V .
Proposition 27. Every (i, j) − SC − R1 is (i, j) − SC − R0 space.
Proof
Suppose that (X, τ1 , τ2 ) is (i, j) − SC − R1 . Let U be (i, j) − SC − open set and
x ∈ U . Then for each y ∈ X − U , (i, j) − SC Cl({x}) 6= (i, j) − SC Cl({y}).
Since (X, τ1 , τ2 ) is (i, j) − SC − R1 , so there exist disjoint (i, j) − SC − open sets
Uy and Vy such that (i, j) − SC Cl({x}) ⊆ Uy and (i, j) − SC Cl({y}) ⊆ Vy . Let
A = ∪{Vy : y ∈ X − U }, then X − U ⊆ A, x ∈ / A and A is (i, j) − SC − open set.
Therefore, (i, j)−SC Cl({x}) ⊆ X −A ⊆ U . Hence, (X, τ1 , τ2 ) is (i, j)−SC −R0 .
The converse of Proposition 27 true, in general, as shown in the following ex-
ample:

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Example 28. Let X = N ∪ {x, y}, where x 6= y and x, y ∈ / N (N : the set of all
natural number). Define τ on X as: A ∈ τ if and only if A ⊆ N or if, x, y ∈ A, then
A contains all N but not finitely many points. Let τ1 = τ2 = τ . Then (X, τ1 , τ2 ) is
(i, j) − SC − R0 but not (i, j) − SC − R1 space.
Corollary 29. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − T2 if and only if it
is (i, j) − SC − R1 and (i, j) − SC − T0 .
Proof Follows from Theorem ?? and Proposition 27, the fact every (i, j)−SC −T2
is (i, j) − SC − T1 [4].
Corollary 30. A bitopological space (X, τ1 , τ2 ) is (i, j) − SC − R1 if and only if
for each x, y in X, (i, j) − SC − ker({x}) 6= (i, j) − SC − ker({y}), there exist
disjoint (i, j) − SC − open sets U and V such that (i, j) − SC Cl({x}) ⊆ U and
(i, j) − SC Cl({y}) ⊆ V .

Acknowledgement

I thank assistant professor Dr. Halgwrd Mohammed Darwesh with (particular

technique or finding example) for comments that greatly improved the manuscript).

References

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spaces”, Amer. Math. Monthly, vol. 68, pp. 886-893, 1961.
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185-189, 1954.
[3 ] Hardi N. Aziz, ”(i, j) − SC − Open set in Bitopological Spaces”, Int. J. of
Scientific and engineering research, 3 (10), 2012.
[4 ] Hardi N. Aziz, Dilan F. Ahmad, Hero M. Salih, ”New Types of Separation
Axioms in Bitopological Spaces”, Salahadin University, zanco journal, vol.
24, no. 4, 2015.
[5 ] Kelly J. C., ”Bitopological Spaces”, Proc. London, Math. Soc. 13, 71-89,
1963.
[6 ] K. K. Dube, ”A note on R0-topological spaces”, Mat. Vesnik, vol. 11(26),
pp. 203-208, 1974.
[7 ] K. K. Dube, ”A note on R1-topological spaces”, Period. Math. Hungar.,
vol. 13, no. 4, pp. 267-271, 1982.
[8 ] Levine N., ”semi-open sets and semi-continuity in topological spaces”,
Amer. Math. Monthly, 70 (1), 36-41, 1963.
[9 ] N. A. Shan in, ”On separation in topological spaces”, C. R. (Doklady)
Acad. Sci. URSS (N.S.), vol. 38, pp. 110-113, 1943.

169
170
 Centralizers of Centrally Semiprime Rings
Adil Kadir Jabbar1 and Farhad Rafiq Krush2
1
Department of Mathematics, School of Science, Faculty of Science and Science
Education, University of Sulaimani, Sulaimani, Kurdistan-Iraq
2
Directorate of Education of Sulaimani- District West
[email protected], [email protected]
Abstract
In this paper, centrally semiprime rings are studied and some conditions are
obtained under which additive mappings of centrally semiprime rings are central-
izers and several properties of centralizers of centrally semiprime rings are proved.
Furthermore, centralizers of centrally semiprime rings are studied and so many
properties of this type of centralizers are also obtained.
Keywords: semiprime rings, centrally semiprime rings, centralizers and Jordan
centralizers, closed Lie ideals and scp-mapping.

Introduction

Let R be a ring. A non-empty subset S of R is called a multiplicative system

in R if 0 ∈ / S and a, b ∈ S implies ab ∈ S. A multiplicative system S in R is
called a central multiplicative system if [S, R]={0}, [S, R]={[s, r] : s ∈ S, r ∈ R}
and [s, r] = sr − rs. It is clear that if S is any central multiplicative system in R,
then S ⊆ Z(R) , where Z(R)={a ∈ R : ar = ra, f orallr ∈ R}. An element a ∈ R
is called a central element if a ∈ Z(R). R is called a semiprime ring if a ∈ R, then
aRa={0} implies that a = 0,where aRa={ara : r ∈ R} and R is called a centrally
semiprime ring if RS is a semiprime ring for each central multiplicative system
S in R. A non zero ideal A of R is called an essential ideal if A ∩ B 6={0}, for
every non zero ideal B of R. A mapping T : R → R is called an additive mapping,
if T (x + y) = T (x) + T (y), for all x, y ∈ R.Let n ∈ Z + , then R is called an
n-torsion free ring if x ∈ R such that nx = 0, then x = 0. A left (resp. a right)
centralizer of R is an additive mapping T : R → R which satisfies T (xy) = T (x)y
(resp. T (xy) = xT (y)), for all x, y ∈ R. An additive mapping T : R → R is
called a Jordan centralizer if T (xy) = xT (y) = yT (x), for all x, y ∈ R, where
xy = xy + yx. An element a ∈ R is called a dependent element of a mapping
F : R → R if,F (x)a = ax holds for all x ∈ R and D(F ) will denote the collection
of all dependent elements of F . An element x ∈ R is called a fixed element of
f : R → R, if f (a) = a.An additive subgroup U of R is called a Lie ideal if
[U, R] ⊆ U and a Lie ideal U of R is called a square closed Lie ideal if u2 ∈ U , for
all u ∈ U . A Lie ideal U of R is called a centrally closed Lie ideal if, U S ⊆ U , for
each central multiplicative system S in R.Let S ⊆ R, then a mapping f : R → R
is called a strong commutativity preserving (simply scp-mapping) on S, if for all
x, y ∈ S, [x, y] = [f (x), f (y)]. Let φ 6= S ⊆ R . An additive mapping f : R → R
is called an anti-homomorphism on S if f (xy) = f (y)f (x), for all x, y ∈ S and f
is said to be an involution on R if the following conditions hold.
1. f (f (a)) = a, for all a ∈ R.
2. f (a + b) = f (a) + f (b), for all a, b ∈ R.

171
 3. f (ab) = f (b)f (a), for all a, b ∈ R.

Additive Mappings of Centrally Semiprime Rings which are

Centralizers

Throughout this section, R will denote a ring (need not contain identity and
need not be commutative) for which Z(R) contains no proper zero divisors, then
one can easily prove that Z(R) \ {0} is a central multiplicatively closed set in R.
Lemma 1. Let R be a ring and S be a central multiplicative system in R and let
d : R → R be a mapping, then d is a centralizer (resp. a left or a right) centralizer
of R if and only if the induced mapping d∗ is a centralizer (resp. a left or a right)
centralizer of RS .
Proof (⇒ ) Let d be a centralizer of R, to show d∗ is a centralizer of RS . By
[5, Lemma3.4], the induced mapping d∗ is an additive mapping. Let as , bt ∈ RS ,
then d∗ (as bt ) = d∗ ((ab)st ) = (d(ab))st = (d(a))b)(st) = (d(a))s bt = d∗ (as )bt . One
can easily show that d∗ (as bt ) = as d∗ (bt ) for all as , bt ∈ RS . Hence, d∗ is a central-
izer of RS .
( ⇐ ) Let d∗ is a centralizer of RS be a centralizer of RS ,so d∗ is additive.
Hence, by [5, Lemma3.4], d is an additive mapping. Let a, b ∈ R , since S 6= φ,
so fix s ∈ S, so that as , bt ∈ RS . Now, (d(ab))ss = d∗ ((ab)ss ) = d∗ (as bs ) =
d∗ (as )bs = (d(a))s (bs )= d(a)b)ss thus (d(ab) − d(a)b)ss= 0 , so by [5, Lemma3.1]
, d(ab) − d(a)b = 0, thus d(ab) = d(a)b. As the same way, one can show
d(ab) = ad(b), for all a, b ∈ R. Hence, d be a centralizer of R.
By the same way we can prove that d is a left (resp. a right) centralizer of R if and
only ifd∗ is a left (resp. a right) centralizer of RS .

Next, in the following results we give some conditions which make an

additive mapping on a centrally semiprime ring as left or as a right centralizer.
Lemma 2. Let R be a 2-torsion free centrally semiprime ring and T : R → R be
an additive mapping that satisfies T (x2 ) = T (x)x , for all x ∈ R , then T is a left
centralizer of R.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and since R is a centrally semiprime ring, so
RS is a semiprime ring. As, R is a 2-torsion free, we can easily get that,RS is also
a 2-torsion free. By [5, Lemma3.4] , T∗ is an additive mapping on RS . Now, for
all as ∈ RS , we have T∗ (a2 s ) = T∗ (as as ) = T∗ ((aa)ss ) = T∗ (a2 ss ) = (T (a2 ))ss =
(T (a)a)ss = (T (a))s as = T∗ (as )as , so by [17, P roposition1.4], one can get that
T∗ is a left centralizer on RS , and by [Lemma2.1], T is a left centralizer on R
Remark 3. It is obvious that, every additive mapping T on a 2-torision free
centrally semiprime ring R which satisfies T (x2 ) = xT (x), for all x ∈ R is a right
centralizer, so that we can state the following lemma.
Lemma 4. Let R be a 2-torsion free centrally semiprime ring. If T : R → R is an
additive mapping which satisfies T (x2 ) = xT (x), for all x ∈ R, then T is a right
centralizer of R.

172
 Proof One can prove it easily by the same technique as in Lemma2.2 âŹę.
Theorem 5. Let m ≥ 1, n ≥ 1 be fixed integers and let R be a (m+n+2)!−torsion
free centrally semiprime ring with an identity element, suppose there exists an
additive mapping T : R → R, such that T (xm+n+1 ) = xm T (x)xn holds for all
x ∈ R, then T is a centralizer of R.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and also RS is a semiprime ring . As, R is a
(m + n + 2)!-torsion free, RS is a (m + n + 2)!-torsion free semiprime ring, and
since T : R → R is an additive mapping, so by [5, Lemma3.4], T∗ : RS → RS
is an additive mapping. Now, by simple computation one can get T∗ (xm+n+1 s )=
(xms )T (x
∗ s )(x n
s ), for all x s ∈ R S . Hence, by [13, T heorem1], T∗ is a centralizer of
RS and by Lemma 2.1, T is a centralizer of R âŹę.
Theorem 6. Let R be a 2-torsion free centrally semiprime ring and let T : R → R
be an additive mapping such that 2T (x2 ) = T (x)x + xT (x) holds for all x ∈ R,
then T is a left as well as a right centralizer of R.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and since R is a centrally semiprime ring, so
RS is a semiprime ring. Also, RS is a 2-torsion free. As, T : R → R is an additive
mapping, by [5, Lemma3.4] , T∗ is an additive mapping on RS . Now for all xs ∈ RS
, where x ∈ R, s ∈ S , then we get 2T∗ (x2 s ) = T∗ (x2 s ) + T∗ (x2 s ) = T∗ (xs xs ) +
T∗ (xs xs ) = T∗ ((x2 )ss ) + T∗ ((x2 )ss ) = (T (x2 ))ss + (T (x2 )ss ) = (T (x2 ) + T (x2 ))ss =
(2T (x2 ))ss )= (T (x)x + xT (x)))ss =(T (x)x)ss + (xT (x))ss = (T (x)x)ss + (xT (x))ss
=(T (x))s xs +xs (T (x))s = T∗ (xs )xs + xs T∗ (xs ). Thus, by[15, T heorem1], T∗ is a
left and a right centralizer of RS and by Lemma2.1, T is a left and a right centralizer
on R âŹę.
Theorem 7. Let R be a 2−torsion free centrally semiprime ring and f be an
involution on R and let T : R → R be an additive mapping such that T (xf (x)) =
T (x)f (x)[T (f (x)x) = f (x)T (x)] holds for all x ∈ R, then T is a left (resp. a right
) centralizer.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and since R is a centrally semiprime ring, so
RS is a semiprime ring and RS is a 2−torsion free and by [5, Lemma3.4], f∗ is an
involution on RS . Since T : R → R is an additive mapping so by[5, Lemma3.4],
T∗ is an additive mapping on RS . Now, for xs ∈ RS , where x ∈ R, s ∈ S , then
T∗ (xs f∗ (xs )) = T∗ (xs (f (x))S = T∗ ((xf (x))ss = (T (x)f (x))ss = (T (x))s (f (x))s =
T∗ (xs )f∗ (xs ) As, the same way one can get T∗ (f∗ (xs )xs ) = f∗ (xs )T∗ (xs ). Hence,
by[14, T heorem1],T∗ is a left (resp. a right ) centralizer on RS , then by Lemma2.1,
T is a left (resp. a right ) centralizer on R âŹę.

At the end of this section, we give some conditions under which a certain
element of a centrally semiprime ring centralizes a given ideal.
Theorem 8. Let R be a ring and I be a non-zero ideal of R and S be a central
multiplicative system in R. If a ∈ R, s ∈ S, then a centralizes I if and only if as
centralizes IS .
Proof Let I be a non-zero ideal of R, so IS is a non-zero ideal of RS .

173
 (⇒ ) Let a centralizes I. To show as centralizes IS . Let xt ∈ IS be any element,
where x ∈ I, t ∈ S. Then, as xt = (ax)st = (xa)ts = xt as .Hence, as centralizes IS
(⇐)The proof of the converse side can be done by the same argument âŹę .
Theorem 9. Let R be a centrally semiprime ring and I be a nonzero ideal of R.
If z ∈ R centralizes the set [I, I], then z centralizes I.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and since R is a centrally semiprime ring, so
RS is a semiprime ring. Since, I is a nonzero ideal of R, then IS is a nonzero ideal
ofRS . Next, for all zv ∈ RS and for all xs , yt ∈ IS , where z, x, y ∈ R, s, t, v ∈ S,
then zv [xs , yt ] = zv (xs yt − yt xs ) = zv (xy)st − ((yx)ts ) = zv ((xy)st − ((yx)st ) =
zv (xy − yx)st = (z(xy − yx))vst = ((xy − yx)z)stv = (xy − yx)st zv = ((xy)st −
(yx)st )zv = (xy)st − (yx)st )zv = (xs yt − yt xt )zv = [xs , yt ]zv , so thatzv centralizes
[IS , IS ], so by [3, Lemma1] , zv centralizes IS , and hence by T heorem2.8 ,z
centralizes I âŹę.

Some Properties of Centralizers of Centrally Semiprime Rings

In below, we prove some properties of centralizers of centrally semiprime

rings and we determine the relation between dependent elements, fixed elements
and central elements of centrally semiprime rings.
Lemma 10. Let R be a centrally semiprime ring in which Z(R) contains no proper
zero divisors of R and T be a left centralizer of R. Then,a ∈ D(T ) if and only if
a ∈ Z(R) and T (a) = a.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0} is
a central multiplicative system in R and since R is a centrally semiprime ring, so
RS is a semiprime ring. By [5, Lemma3.4], the induced mapping T∗ is an additive
mapping, and since T is a left centralizer of R, so by Lemma2.1, T∗ is a left
centralizer of RS .
(⇒) Let a ∈ D(T ) and fix ans ∈ S , then by[5, Lemma3.5],as ∈ D(T∗ ) and by
[12, T heorem2.1],as ∈ Z(RS ) and T∗ (as ) = as and as Z(R) contains no proper
zero divisors of R one can easily get that a ∈ Z(R) and T (a) = a.
(⇐) Let a ∈ Z(R) and T (a) = a. So, for any s ∈ S, then as ∈ (Z(R))S ⊆ Z(RS )
and (T (a))s = as and thus T∗ (as ) = as . So again by [12, T heorem2.1],as ∈ D(T∗ )
, and by[5, Lemma3.5], a ∈ D(T ) âŹę.
Lemma 11. Let R be a centrally semiprime ring with identity 1 in which every non
zero element of Z(R) is a unit in R. If T is a left centralizer of R and a ∈ D(T ),
then T is identity on the ideal J of R generated by a. (Note that the notation < a >
is used for the ideal of R generated by a ).
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R)subminuse {0}
is a central multiplicative system in R and since R is a centrally semiprime ring,
so RS is a semiprime ring and clearly, 1 ∈ S. As, T is an additive mapping,
so by [5, Lemma3.4], the induced mapping T∗ is an additive mapping and by
Lemma2.1, T∗ is a left centralizer of RS . Fix an s ∈ S and since a ∈ D(T )
, so by [5, Lemma3.5], as ∈ D(T∗ ). So,T∗ is an identity on the ideal K of RS
generated by as and then by [6, P roposition2.16] , there exists an ideal J of R such

174
that JS = K =< as >.To show J =< a >. Now,as ∈ JS , so ta ∈ J , for some
t ∈ S = Z(R) \ {0}, hence, t−1 ∈ R and then a = t−1 ta ∈ J , so that < a >⊆ J.
Next, Let x ∈ J, then x1 ∈ JS =< as > , so that x1 = rt as + as bd , for some
rt , bd ∈ RS . Then, x1 = rt as + as bd
x1 = (rad)tsd + (abt)sdt = (rad + abt)sdt . Hence there exists q ∈ S such that
q(tsdx − (rad + abt)) = 0, then by [5, Lemma3.1], tsdx = rad + abt and as
t, s, d ∈ S ,then tsd ∈ S = Z(R) \ {0}, and hence (tsd)−1 inR, so that x =
s−1 t−1 ra + as−1 d−1 b ∈< a > . So that, J ⊆< a >. Hence, J =< a >. To show
T is identity on the ideal of R generated by a. Now, if x ∈ J =< a >, then for
any s ∈ S, so xs ∈ JS =< as > and hence T∗ (xs ) = xs , so ãĂŰ(T (x))s = xs ,
thus (T (x))s − xs = 0 , then (T (x) − x)s = 0, so by [5, Lemma3.1], T (x) − x = 0,
that is, T (x) = x, where x ∈ J =< a >. Thus, T is an identity on the ideal J
generated by a âŹę.
Now, we prove that under certain conditions the set of all dependent elements
of a mapping on a ring forms an ideal if and only if the set of all dependent elements
of the induced mapping forms an ideal of the localized ring.
Lemma 12. Let R be a ring in which Z(R) contains no proper zero divisors of R
and S is a central multiplicative system in R. If T : R → R is a mapping, then
D(T ) is an ideal of R if and only if D(T∗ ) is an ideal of RS .
Proof (⇒) Let D(T ) be an ideal of R. Since, D(T ) 6= φ, so there exists
a ∈ D(T ). Let s ∈ S, then by [5, Lemma3.5], as ∈ D(T∗ ), that is φ 6= D(T∗ ) ⊆ RS .
Let, cu , bt ∈ D(T∗ ), xv ∈ RS , where c, b, x ∈ R, u, t, v ∈ S, then by [5, Lemma3.5],
c, b ∈ D(T ). Now, cu − bt = (ct − bu)ut and since, ct − bu ∈ D(T ), so by
[5, Lemma3.5], (ct − bu)ut inD(T∗ ) , that is cu − bt ∈ D(T∗ ). Also, cu xv = (cx)uv
and since cx ∈ D(T ), so by [5, Lemma3.5], (cx)uv ∈ D(T∗ ), thus cu , xv ∈ D(T∗ ) .
As, the same way one can get xv cu ∈ D(T∗ ). So, D(T∗ ) is an ideal of RS .
(⇐) Let D(T∗ ) be an ideal of RS , to show D(T ) is an ideal of R. Since, D(T∗ ) 6= φ,
so there exists ds ∈ D(T∗ ), for some d ∈ R and s ∈ S and by [5, Lemma3.5],
d ∈ D(T ), that is D(T ) 6= φ , and thus φ 6= D(T ) ⊆ R . Let, y, b ∈ D(T ), x ∈
R. If s ∈ S , then ys bs ∈ D(T∗ ) and xs ∈ RS . Now, (y − b)s = ys − bs ∈
(T∗ ), so by [5, Lemma3.5] ,y − b ∈ D(T ). Also, (yx)ss = ys xs ∈ D(T∗ ), so by
[5, Lemma3.5],yx ∈ D(T ). As the same way one can get, (xy)ss = xs ys ∈ D(T∗ )
, that means xy ∈ D(T ), so that D(T ) is an ideal of R âŹę.
Next, we prove some lemmas in which we give some conditions which make the
set of all dependent elements of certain mappings of centrally semiprime rings as
subrings or as ideals.
Lemma 13. Let R be a centrally semiprime ring with identity in which every non-
zero element of Z(R) is a unit in R and T is a left centralizer of R, then D(T ) is
a semiprime subring of Z(R).
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0}
is a central multiplicative system in R and since R is a centrally semiprime ring,
so RS is a semiprime ring. By [5, Lemma3.4] ,T∗ is an additive mapping and
since T is a left centralizer of R, so by Lemma2.1, T∗ is a left centralizer of
RS so, by [12, T heorem2.3], D(T∗ ) is a semiprime subring of Z(RS ). Now, to
show D(T ) is subring of Z(R). Since,D(T∗ ) 6= φ, so, there exists as ∈ D(T∗ ),
for a ∈ R, s ∈ S and by [5, Lemma3.5], a ∈ D(T ), that is D(T ) 6= φ. Now,let

175
 x ∈ D(T ), so by [5, Lemma3.5], xs ∈ D(T∗ ) ⊆ Z(RS ). If u ∈ R, then us ∈ RS
and hence xs us = us xs , thus (xu)ss = (ux)ss , then (xu − ux)ss = 0 and by
[5, Lemma3.1], xu − ux = 0, thus xu = ux, for all u ∈ R , so that x ∈ Z(R) and
thus D(T ) ⊆ Z(R), thus φ 6= D(T ) ⊆ Z(R). Also, it is easy to show that, for all
a, b ∈ D(T ), then a − b ∈ D(T ) and ab ∈ D(T ). Hence, D(T ) is a subring of
Z(R). Now, let xD(T )x = {0}. If av ∈ D(T∗ ), for a ∈ R, v ∈ S , then a ∈ D(T ),
so that xs av xs = (xax)svs = 0, so that xs D(T∗ )xs = {0}, thusxs = 0, from which
one getsx = 0, so that D(T ) is a semiprime subring of Z(R) âŹę.
Corollary 14. Let R be a commutative centrally semiprime ring in which Z(R)
contains no proper zero divisors of R and T is a centralizer of R, then D(T ) is an
ideal of R.
Proof Since Z(R) contains no proper zero divisors of R, so S = Z(R) \ {0}
is a central multiplicative system in R and since R is a centrally semiprime ring,
so RS is a semiprime ring and as R is commutative,RS is commutative and by
Lemma2.1, T∗ is a centralizer of RS , so by [12, Corollary2.4], D(T∗ ) is an ideal
of RS and by Lemma3.3, D(T ) is an ideal of R âŹę.
In the remaining results of this section we prove that under certain conditions one
can transform the essentially property of ideals of a given ring to the localized ring.
Lemma 15. Let R be a ring with identity in which every non zero element of Z(R) is
a unit in R and U, V be ideals of R. If S is a central multiplicative system in R, then.

1. U + V is an ideal of R if and only if US + VS is an ideal of RS .

2. U ∩ V = {0} if and only if US ∩ VS = {0}.
3. U + V is an essential ideal of R if and only if US + VS is an essential ideal
of RS .
Proof
[1] The proof is straightforward.
[2] Let U ∩ V = {0}. If US ∩ VS 6= {0}, then there exists 0 6= as inU S ∩ VS ,
this implies that as ∈ US and as ∈ VS . It is easy to prove that a ∈ U and a ∈ V ,
so a ∈ U ∩ V and as as 6= 0, we have a 6= 0, which is a contradiction, so that
US ∩ VS = {0} .
Conversely, let US ∩ VS = {0} . If U ∩ V 6= {0} , then there exists 0 6= a ∈ U ∩ V .
Now, fix an s ∈ S , then as ∈ US and as ∈ VS and then ãĂŰas ∈ US ∩ VS and
so that as = 0 and by [5, Lemma3.1], a = 0, which is a contradiction. Hence,
U ∩ V = {0}.
[3] Let US ∩ VS be an essential ideal of RS .We will show that U + V is an
essential ideal of R. Clearly, U + V is an ideal of R. Let J be an ideal of R
such that (U + V ) ∩ J = {0}, then JS is an ideal of RS , and (JS + VS ) ∩ JS =
((U + V ) ∩ J)S = 0, so that JS = 0. Hence, we get J = 0. Thus, U + V is an
essential ideal of R.
Conversely, let U +V be an essential ideal of R. Let K be any ideal of RS such that
(US + VS ) ∩ K = {0} , then by [6, P roposition2.16], there exists an ideal J of R
such that K = JS , then (US + VS ) ∩ JS = {0}, so that ãĂŰ((U + V ) ∩ J)S = {0}.
By [5, Lemma3.1], we get that (U + V ) ∩ J = {0}, which gives that J = {0}, thus
JS = {0}, so that US + VS is an essential ideal of RS âŹę.

176
Corollary 16. Let R be a ring with identity in which every non zero element of
Z(R) is a unit in R and T be a centralizer of LR. If S is a central multiplicative
system in RLand U, V are ideals of R, then U V is an essential ideal of R if and
only if US VS is an essential ideal of RS .
Proof The proof follows directly from Lemma3.6, and from the fact that U ∩V =
{0} if and only if US ∩ VS = {0} âŹę.
Lemma 17. Let R be a ring with identity in which every non zero element of Z(R)
is a unit in R and T is a centralizer of R. If U is an ideal of R and S is a central
multiplicative system in R, then.
1. T = I on U if and only if T∗ = I∗ on US , where I∗ is the induced map on
RS .
2. D(T |U ) = {0} if and only if D(T∗ |(US ) ) = {0}.
Proof
[1] (⇒) Let T = I on U and to showT∗ = I∗ on US . Let, xt ∈ US , where
x ∈ U, t ∈ S, then T∗ (xt ) = (T (x))t = (I(x))t = I∗ (xt ), thus T∗ = I∗ .
(⇐) Let, T∗ = I∗ on US . Let,y ∈ U . Since, S 6= 0 , so fix an s ∈ S, then
(T (y) − y)s = (T (y))s − ys = T∗ (ys ) − ys = 0, thus by [5, Lemma3.1], we get that
T (y) − y = 0, so that T (y) = y = I(y). Hence,T = I on U .
[2](⇒) Let,D(T |U ) = {0}. Let, at ∈ D(T∗ |US , then by [5, Lemma3.5], a ∈ D(T |U ),
so that a = 0, thus at = 0. Hence, D(T∗ |US )) = {0}.
(⇐) Let, D(T∗ |US ) = {0}. By the same technique as in the first part we get
D(T |U ) = {0} âŹę.
Lemma 18. Let R be a commutative centrally semiprime ring with identity in which
every non zero element of Z(R) is a unit in R and T is a centralizer of R. Then,
there are ideals U and V of R such that

L
1. U V is an essential ideal of R.
2. T = I on U , T (V ) ⊆ V .
3. D(T |V ) = {0} where T |V is the restriction of T on V
Proof Since every non-zero element of Z(R) is a unit in R, so S = Z(R)\{0} is a
central multiplicative system in R and RS is a semiprime ring. By [5, Lemma3.4],
T∗ is an additive mapping and since R is commutative, so RS is commutative.
Since T is a centralizer of R, so by Lemma2.1,T∗ is a centralizer of RS . Hence,
by [12, T heorem2.6], there exist ideals K and M of RS such that
L
1. K M is an essential ideal of RS .
2. T∗ = I∗ on K ,T∗ (M ) ⊆ M .
3. D(T∗ |M ) = {0} , where T |M is the restriction of T∗ on M .
By [6, P roposition2.16],
L L exist ideals U, V L
there of R such that K = US , M =
VS , that is K M = US LVS and since US VS is an essential ideal of
RS , so by Corollary3.7,U V is an essential ideal of R, which proves (1)
. Since, T∗ = I∗ on K, by Lemma3.8, T = I on U . We have T∗ (M ) ⊆ M
, that means T∗ (VS ) ⊆ S. To show, T (V ) ⊆ V , let y ∈ T (V ), then there

177
 exists x ∈ V such that y = T (x). If s ∈ S is any fixed element in S,
then ys = (T (x))s = T∗ (xs ) ∈ VS , so there exists u ∈ V, t ∈ S such that
ys = ut , then there exists r ∈ S such that r(ty − su) = 0, from which we get
y = t−1 su ∈ V . Hence, T (V ) ⊆ V and this proves (2).
Now, D(T∗ |M ) = {0},that is, D(T∗ |(VS ) ) = {0}, so by Lemma3.8, we get
D(T |V ) = {0}, which proves (3) âŹę.

References
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on semiprime rings. Bull. Korean Math. Soc. 43 (2006), No. 4, pp. 711-713.
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rings with involution, International Journal of Algebra, Vol. 1, No. 12, (2007),
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[3 ] M. N. Daif and H. E. Bell: Remarks on derivations on semiprime rings, Internat.
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152, (2007), pp. 135-150.
[5 ] Jabbar, A. K. and Krush, F. R. : A generalization of Semiprime rings, Journal
of Garmian University, Vol. 1, 2015, pp34-47.
[6 ] Jabbar, A. K.: A generalization of prime and weakly prime submodules, Pure
Mathematical Sciences, vol. 2, 2013, no. 1, 1-11.
[7 ] A. K. Jabbar and A. H. Majeed: On Centrally Semiprime Rings and Cen-
trally Semiprime Near-rings with Derivations, Journal of Kirkuk University,
Scientific Studies vol. 3, No.1, (2008), pp.158-168.
[8 ] A. K. Jabbar and A. H. Majeed: some properties of centrally closed Lie ideal
of centrally prime rings, Journal of Al-Nahrain University, Vol. 10 (2), (2007),
pp. 162- 167.
[9 ] Y. S. Jung and K. H. Park: On Generalized (Μ,Κ)- Derivations and com-
mutativity in Prime Rings, Bull. Korean Math. Soc., Vol. 43, (2006), pp.
101-106.
[10 ] M. D. Larsen and P. J. McCarthy: Multiplicative Theory of Ideals, Academic
Press New York and London (1971).
[11 ] M. Samman and N. Alyamani: Derivations and Reverse Derivations in Semiprime
Rings, International Math. Fourm, No. 39, (2007), pp.1895-1902.
[12 ] M. S. Samman and M. A. Chaudhry: Dependent elements of left centralizer
of semiprime rings, The Arabian Journal for Science and Engineering, vol. 33,
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[13 ] I. K. UIBI: A remark on centralizers in semiprime rings, Glasnik Matematicki,
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[14 ] J. Vukman and I. K. UIBI: On centralizers of semiprime rings with involution,
Studia scientiarum Mathematicarum Hungarica, 43(1), (2006), pp. 61-67.
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Math. Univ. Carolin. 40,3(1999), pp. 447-456.
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line 32,4, (1991), pp. 609-614.

179
180
Computer Session
Design Security System Based on AES and MD5 for Smart Card
Raghad Z.Yousif1 , Shahab.W.Kareem2 and Ammar. O.Hasan3
1
College of science, Salahaddin University- Erbil, Kurdistan-Iraqi
2
Erbil Technical Engineering College-Erbil Polytechnic University, Kurdistan-Iraqi
3
College of science, Salahaddin University- Erbil, Kurdistan-Iraqi
[email protected], [email protected] and [email protected]
Abstract
In the most and variety of applications, it is necessary to build a security system.
In this paper a design and implement of secure system has been presented. This
system consists of memory card, memory card reader, AES algorithm for data
security (encryption and decryption data) and MD5 algorithm for authentication and
data transection. The memory card is a serial EEPROM (AT24CXX) connected
with the memory card reader by USB connector. This memory is used to store
the personal information such as identification data. The memory card reader is
the second part of the system which contains the AT89C51 microcontroller, clock,
reset, and power circuits. The card reader is used to generate power and control
signals also for transferring data between PC and memory card through serial bus.
Software implementation and testing of proposed system which were written in
ModelSim6.2. the AES algorithm used to perform the encryption and decryption
data and then authentication from this data using MD5 algorithm has been used
to recognize the patterns which are stored in the serial EEPROM . MD5 provides
the capability to store many associated pattern pairs and fast training while AES
symmetric key cryptography algorithm has been used for programming the serial
EEPROM to prevent data from external intrusion.
Keywords: Advanced Encryption Standard (AES), MD5 message-digest algorithm,
Electrically Erasable Programmable Rea- Only Memory (EEPROM), memory card,
memory card reader.

Introduction

One of the most security mechanisms is detect and prevent attacks. The requirement
of system security is analyzing the system, understanding the mechanisms of secu-
rity policy and knowledge of the related assumptions and trust, which lead to the
threats and the degree to which they may be realized. Such knowledge allows one
to design better mechanisms and policies to neutralize these threats. This operation
leads to risk analysis. Computer security based on confidentiality, integrity and
availability. The clarification of these three aspects vary, as to the contexts in which
they arise. The clarification of an aspect in a given environment is dictated by the
needs of the individuals, customs, and laws of the particular organization [1]. The
security policyŠs specification to ŞsecureŤ and Şnon-secureŤ actions, these security
mechanisms can prevent the attack, detect the attack, or recover from the attack.
The strategies may be used together or separately. An attack will fail is called
Prevention. Preventative mechanisms often are very cumbersome and interfere with
system use to the point that they hinder normal use of the system. Detection is most

181
 useful when an attack cannot be prevented, but it can also indicate the effectiveness
of preventative measures. Detection mechanisms accept that an attack will occur;
the goal is to determine that an attack is under way, or has occurred, and reports it.
The attack may be monitored, however, to provide data about its nature, severity,
and results. Typical detection mechanisms monitor various aspects of the system,
looking for actions or information indicating an attack [2].Smart cards are often used
in different applications that needed more strong security protection and authenti-
cation, such as banking, transportation, and identification. In these applications,
major focuses are on the authentication and secure communication links between
cards readers [3].The objectives of this paper are First to design and implement a
flexible digital security system. The system can easily change its function through
the design of multifunction memory card that can be used in different fields. Second
Using a serial EEPROM as a protected memory card through designing MATLAB
Graphic User Interface programs (GUI) to implementing AES cryptography and
MD5algorithms used for store and retrieve the protected data from memory card
based on AT89c51 microcontroller that used to connect memory card with PC. The
remainder of this paper is organized as follows: In Section 2, introduction of cryp-
tographic and smart card. Section 3 presents types, architectures and programming
for microcontrollers.in section 4 System Design and Implementation, Finally, the
last section provides concluding remarks and directions for future work.

Secure System

In this section discussed the security implemented in this system.

Advanced Encryption Standard (AES)

AES is a symmetric key encryption algorithm for cryptography, AES based on the
Rijndael block cipher by Joan Daemen and Vincent was selected by the National
Institute of Standards and Technology (NIST) to become the AES as a replacement
for the outdated Data Encryption Standard (DES) algorithm [4]. It is a substitu-
tion permutation network (SPN) depended on iterative block cipher which acts on
plaintext blocks of 128-bits and supports significantly larger key sizes than DES,
i.e. 128-bits, 192-bits or 256-bits. Based on the key size, the number of rounds
is either 10, 12 or 14 respectively [5]. For the work here, the 256-bit key size is
examined, although only power traces from the first round are utilized hence the
attack is equally viable for the smaller key sizes. Figure 1 outlines a pseudo code
description of the AES algorithm. First, the plaintext block B is copied into the
state variable, which is a 4 ¤ 4 matrix of bytes. Then, an start AddRoundKey
function simply XORs the initial key to the state. This is followed by identical
round transformations consisting of the functions; S-Box, ShiftRows, MixColumn,
and AddRoundKey. The final round skips the MixColumn operation to generate the
ciphertext. [6][7][8].

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 Figure 1: Advanced Encryption Standard procedure

Message digest 5

Hash function is used to warranty the originality and authentication of the text of
file. Encryption does not protect data from changing by another party. Need a way
to guarantee that data received at destination in its original form as sent by the
transmitter and it is coming from an authenticated source simply it should be in the
unique manner. When hash function is applied to the message or the original data
to send, it always produces fixed length value. Message digest 5 or MD5 [9] [10] is
one of the most message-digest algorithm that was sophisticated by Ron Rivest in
1991. MD5 algorithm takes an input in different length and produces an output in
the form of a digest with the fixed length of 128 bits. Input arrived by this algorithm
will be processed in a block size of 512 bits, which will then be divided into 16
sub blocks, each is 32 bits.. The advantages of MD5 [11] [12] is that it is simple to
compute the hash value for any given message, infeasible to find a message that has
a given hash, infeasible to modify a message without changing its hash, impossible
to find two different messages with the same hash. It requires only less memory and
low processing time. All these advantages make it acceptable over other hashing
algorithms and to be used in a security platform. Also MD5 is suitable for detect
the errors that may occur during the transmission. Figure 2 shows the processing
of input on MD5 algorithm.

Smart-Card

Smart cards are the youngest members of the plastic card family. Smart card is
defined as: A plastic card, usually similar in size and shape to a credit card, con-
taining a microprocessor and memory (which allows it to store and process data)
and complying with ISO 7816 standard [14]. Smart cards can be classified on the
basis of various parameters, this paper present smart-card classification on basis of

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 Figure 2: The procedure of the input at 512-bit block of the MD5 [13]

card components, and card interface, as shown in Figure 3.

Figure 3: Smart Cards Classification component

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 Smart cards can be divided into two categories. Chip card that have a processor or
microprocessor cards (EEPROM, ROM, CPU and RAM), and memory card without
a chip that contain (EEPROM, ROM, Security Logic, Straight Memory Cards and
Stored Value Memory Cards) [15].Certification application smart card chip in there
with the user identity-related data, by the legal users to carry, you must log into a
dedicated IC card reader to read the information in order to check the user’s identity.
IC card authentication is based on ”what you have” means. Through the IC card
can guarantee that the hardware cannot copy the user’s identity will not be fake [14]
[15][16].

Microcontroller

A microcontroller is a single chip computer, or is used to describe a system that

includes a minimum of a microprocessor, program memory, data memory, and
inputŰoutput (I/O) [17].

8051 microcontroller

The 8051 microcontroller widely used in security applications, because of its com-
pact and easy implementation, is subjected to easy power analysis attack. Many
controller architectures recently developed have just widened the functional capac-
ity of 8051 controller. Since the underlying technology remains the same of the
recent structure’s vulnerability to power analysis attack is also proved with this ex-
periment. In order to acquire the device foot print from its power consumption, a
high accuracy data acquisition system is required [18] the microcontrollers have an
8-bit data bus. They are capable of addressing 64K byte of program memory and a
separate 64K byte of data memory. The 8051 has 4K byte of code memory imple-
mented as on-chip Read Only Memory (ROM). The 8051 has 128 bytes of internal
Random Access Memory (RAM). The 8051 has two timer/counters, a serial port,
4 general purpose parallel input/output ports, and interrupt control logic with five
sources of interrupts. Besides internal RAM, the 8051 has various Special Function
Registers (SFR), which are the control and data registers for on-chip facilities Pro-
gramming the various internal hardware facilities of the 8051 is achieved by placing
the appropriate control words into the corresponding SFRs [18][19]. Figure 4 show
The 8051 Microcontroller Architecture.

AT89C51 Microcontroller

The AT89C51 is a low power, high-performance CMOS 8-bit microcomputer with

4Kbytes of Flash programmable and erasable read only memory (EEROM). The
device is manufactured using AtmelŠs high-density nonvolatile memory technology
and is compatible with the industry-standard MCS-51 instruction set and pinout
[20]. The on-chip Flash allows the program memory to be reprogrammed in-system
or by a conventional nonvolatile memory programmer. By combining a versatile
8-bit CPU with Flash on a monolithic chip, the Atmel AT89C51 is a powerful

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 Figure 4: The 8051 Microcontroller Architecture.

microcomputer which provides a highly-flexible and cost-effective solution to many

embedded control applications [21]. In this work the AT89C51 microcontroller
using as a main part in card reader to perform transferring data between card reader
and PC. All Atmel Flash microcontrollers have separate address space for Program
and Data Memory, as shown in Figure 5.a. The logical separation of Program and
Data Memory allows the Data Memory to be accessed by 8-bit addresses, which
can be more quickly stored and manipulated by an 8-bit CPU. Nevertheless, 16-
bit Data Memory addresses can also be generated through the DPTR register [20].
Program Memory can only be read, not written to. There can be up to 64K bytes
of Program Memory. The AT89C51 can directly address up to 64K bytes of data
memory external to the chip. The AT89C51 has 128 bytes of on-chip RAM plus
a number of Special Function Registers (SFRs). The lower 128 bytes of RAM can
be accessed either by direct addressing or by indirect addressing. Figure 5.b shows
the AT89C51 data memory organization [22].

To program a microcontroller needed two basic tools. One is the hardware-

programmer board and the second is the software for downloading the program.
The hardware programmer board is used to hole the chip and enable the computer to
download the program into the micro. The programmer board is serially connected
to the computer. The software will need to program the micro will be two fold.

System Design and Implementation

The practical part of the work consists of the stages of building the system hardware
and software.

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 (a) AT89C51 Program Memory (b) AT89C51 Data Memory

Figure 5: Memories of AT89C51

Hardware System

The system hardware consists of two separated parts one for memory card reader
and other for memory card, while the system software consists from all the programs
that are written in personal computer and microcontroller. The Smart Card Based on
AES and MD5 Data Security System hardware as shown in Figure 6 and Figure 7.
The AT89C51 microcontroller is used as a main part in the system. Some additional

Figure 6: Block Diagram of the System Hardware and Software

circuits are connected to this unit to perform the correct operation such as crystal,
reset circuit, RS-232C Interface, Power Supply and Serial EEPROMS. A serial
EEPROM is created with the same technology used in larger parallel EEPROMs. In
this application using the AT24C04 is used because this type of the serial EEPROM
is most popular and inexpensive. The Bidirectional Data Transfer Protocol utilized
by the AT24CXX family allows a number of compatible devices to share a common
2-wire bus. The bus consists of a serial clock (SCL) line and a serial data (SDA)
line. The clock is generated by the bus master and the data is transmitted serially
on the data line, and the most significant bit is first, synchronized to the clock. The

187
 Figure 7: Hardware System

protocol supports bidirectional data transfers in 8-bit bytes. The microcontroller

serves as the bus master, initiating all data transfers and generating the clock which
regulates the flow of data. The serial devices present on the bus are considered
slaves, accepting or sending data in response to orders from the master. The bus
master initiates a data transfer by generating a start condition on the bus. The SDA
pin is normally pulled high with an external device. Data on the SDA pin may
change only during SCL low time periods (referring to Data Validity timing as in
figure 8.a.All addresses and data words are serially transmitted to and from the
EEPROM in 8-bit words. The EEPROM sends a zero to acknowledge that it has
received each word. This happens during the ninth clock cycle (referring to Output
Acknowledge timing as in figure 8.b.

(a) Output Acknowledge timing (b) Data Validity timing

Figure 8: Timing signals

System Software

The system software consists of microcontroller software and PC software. The

microcontroller software contains three phases. The First Phase is configuration

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process for all parameters used inside the algorithms necessary for data transfer to
and from the serial EEPROM. These parameters include the creation of registers,
counters, addresses, commandsĚetc. Additional to the microcontroller serial port
configuration, it is used in data transfer (Data bits= 8bits, Parity = none, Stop bit
= 1bit, Baud rate =19200MHZ). Second Phase Read and Transmit Data is most
important phase written inside the microcontroller and through which we can read
the contents of the serial EEPROM serially by linking the P1.0 pin and P1.1 pin of
microcontroller to the SCL pin and SDA pin for the serial EEPROM, respectively.
Starting to generate the timing signal necessary for the reading process from memory
through the SDA pin of the serial EEPROM, it places the data in the accumulator
register. After storing data in the accumulator register, the microcontroller send, it
to the personal computer through the serial port and then re-processes for the read
and transmit of other data. The Last Phase is Receive and Write Data: This phase
is used for the purpose of changing the contents of the serial EEPROM through
calling the subroutine that withdraw the data from the computer and places it in the
accumulator register and then sends it to the serial EEPROM through the SDA pin
after the generation of signals necessary for the process of writing. PC Software
is part of the software is written by the MATLAB program, which supports the
possibility of designing graphic user interface (GUI). This program is designed to
give different orders of the

system such as sending and receiving data to and from the microcontroller,
encrypting data transmitted, and decryption of data receiving, and sending the data
receive from the serial port to the neural network to identify it. Figure( 9) shows
the main object of the GUI, START button is to start the operations of the system
and initialization of all parameters used in the system. READ FROM SERIAL
EEPROM This button is used to create a serial port object to withdraw data from
the microcontroller to the PC and send it to the security system. In this section the
created serial port must have the same values used in the microcontroller serial port
to ensure that the data is not lost during the transfer process. DECRYPTION button
the data taken from the SERIAL EEPROM by read operation is to be decrypted by
Decryption pushbutton and processed by AES based on the Rijndael block cipher
algorithm. CONVERT TO BINARY VECTOR is taken from the DECRYPTION
callback and converted it to binary vector format. DATA PROCESSING is one of
the most important parts of the system where we use the MD5 for the purpose of
processing and discriminating the data coming from the serial EEPROM. DATA
TO BE WRITTEN is withdraws the data from the STATIC TEXT and convert it
into the sequence numbers and then sends it to the ENCRYPTION part, in addition
to computing the number of bytes that will be sent to the microcontroller through
WRITE TO SERIAL EEPROM callback. ENCRYPTION In this part, the cipher text
of the data to be sent is to be determined by the same algorithm used in decryption
process (AES based on the Rijndael block cipher algorithm).WRITE TO SERIAL
EEPROM In this section serial port object will be generated by the same method
used in read operation, to employ it for transferring encryption data from PC to
microcontroller.

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 Figure 9: Proposed system GUI

System Operation

The operation of the system can be done by the following two main parts, Program-
ming Memory Card and Retrieve Data from Memory Card. Programming Memory
Card execute the following steps using sequentially for the purpose of writing in the
serial EEOROM:
1. Press the START pushbutton to start the operation of the system and initializing
all parameters used in the Graphic User Interface (GUI) to perform the RSA
cryptography algorithm and BAM neural network algorithm.
2. Input the data to be sent to the serial EEPROM through STATIC TEXT and
then press on the DATA TO BE WRITTEN pushbutton.
3. Encrypting the data to be sent through pressing ENCRYPTION pushbutton.
4. Sending the data to the memory card through microcontroller by pressing
WRITE TO SERIAL EEPROM pushbutton. Noting the red LED for storing
Data in the card memory.
While the purpose of Retrieve Data from Memory Card is reading the data from
the memory card, apply the following steps sequentially:
1. Read the contents of the serial port through pressing READ FROM SERIAL
EEPROM pushbutton. Noting the LED for loading Data to the PC.
2. Decrypting the data through pressing DECRYPTION pushbutton.

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 3. Convert the data to sequence number format using CONVERT TO BINARY
VECTOR pushbutton.

4. Insert data input to the neural network for the purpose of processing DATA
PROCESSING pushbutton.

To test the system performance the ModelSim 6.2 has been used .Complete
system model has been written in VHDL. A pattern of size 4Kbyte is prepared, and
simulated in the next step encryption the pattern by AES algorithm (convert it to
ciphertext) as shown in figure (10), while in figure (11) the corresponding signal
pattern in ModelSim for the encryption process has been depicted.

Figure 10: Convert Plain Text to Cipher Text

Figure 11: Modalism simulation results for encryption (Plain Text to Cipher Text)

After encryption pattern, send it to the memory card and verify from it by
programmer. After than withdraw the data from memory card to the PC through
microcontroller and present it to the AES algorithm to convert it to the plaintext for
the purpose of comparing it with the original pattern in step(1) and verify from the
AES algorithm operation is correct, as shown in Figure(12). Figure (13) shows the
ModelSim signaling result of Decryption process in proposed system.

191
 Figure 12: Convert Cipher Text to Plain Text

Figure 13: Modalism simulation results for Decryption (Cipher Text to Plain Text)

Conclusion

The design of the security system is flexible and can easily be improved by replacing
its hardware elements with others of high quality, such as large memory capacity
so that it can be used in various applications. Using MD5 to provide the capability
to store and retrieve information based on its content rather than storage address.
The important characteristic of MD5 is that for an input stimulus that obtains a
specific response, we will get the same response through association if similar input
is presented at a different time. Another important advantage of this technique is
related with the free-of-error match that it gives when tested with other trained
patterns. The MD5 is stronger than other techniques. The AES symmetric key
cryptosystem is the most popular algorithms used in wide range of application.

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