DEPARTMENT OF MATHEMATICS

MATHEMATICS FOR ENGINEERS

I/IV, B.Tech – II Semester
A.Y.2019-2020
COURSE MATERIAL

Course Code : 19MT1201

L-T-P-S structure : 3-0-0-0
Credits : 3
COURSE OUTCOME 1
Introduction to Matrix Algebra
Session 1

LU- Decomposition Method

After reading this chapter, you should be able to:
1. identify when LU decomposition is numerically more efficient than Gaussian elimination,
2. decompose a nonsingular matrix into LU, and
3. show how LU decomposition is used to find the inverse of a matrix.
Like Gauss elimination, LU decomposition method is a kind of exact solution of system of linear
algebraic equations. This method attempts to decompose coefficient matrix into two lower and
upper triangular matrices.
Let A be a square matrix. An LU factorization refers to the factorization of A, with proper row and
or column orderings or permutations, into two factors, a lower triangular matrix L and an upper
triangular matrix U,

In the lower triangular matrix all elements above the diagonal are zero, in the upper triangular
matrix, all the elements below the diagonal are zero. For example, for a 3-by-3 matrix A, its LU
decomposition looks like this:

This method is that any matrix A can be expressed as the product of a lower triangular matrix and
an upper triangular matrix provided all the principal minors of A are non-singular.
If and minors are non-singular i.e.

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Then all these factors if exists is unique.
Let us now consider the equations
a11 x + a12 y + a13 z = b1
a21 x + a22 y + a23 z = b2
a31 x + a32 y + a33 z = b3.
This system of equations can also be written as

Now, let us suppose

Now (1) can be written

as (2)

Put, (3

Equation (2) becomes , which is equivalent to

This led to the solution of and hence V becomes available. Equation (3) becomes

From the above system of equations become available by back substitution.

Computation of L and U matrices

As
Now comparing the left hand multiplication with right hand matrix, we get
Multiplication of first row of L with U and thereafter comparison gives

Multiplication of first column of U with L and thereafter comparison gives

2
Multiplication of second row of L with U and thereafter comparison gives

Multiplication of second column of U with L and thereafter comparison gives

Multiplication of third row of L with U and thereafter comparison gives

Proceeding similar way for higher order matrices and generalizing the method, results in easily
computation of matrices L and U.
Example
Solve the following equations by using LU decomposition method
3x+2y+7z = 4; 2x+3y+z = 5; 3x+4y+z = 7
1 0 0 𝑢 𝑢 𝑢 3 2 7
Let 𝑙 1 0 0 𝑢 𝑢 = 2 3 1
𝑙 𝑙 1 0 0 𝑢 3 4 1
𝑢 𝑢 𝑢 3 2 7
𝑙 𝑢 𝑙 𝑢 +𝑢 𝑙 𝑢 +𝑢 = 2 3 1
𝑙 𝑢 𝑙 𝑢 +𝑙 𝑢 𝑙 𝑢 +𝑙 𝑢 +𝑢 3 4 1
𝑢 =3 𝑢 =2 𝑢 =7
𝑙 𝑢 =2 𝑙 𝑢 +𝑢 =3 𝑙 𝑢 +𝑢 =1
𝑙 𝑢 =3 𝑙 𝑢 +𝑙 𝑢 =4 𝑙 𝑢 +𝑙 𝑢 +𝑢 =1
𝑙 = 2/3 𝑙 =1 𝑢 = 5/3 𝑢 = -11/3 𝑙 = 6/5 𝑢 = -8/5

1 0 0 3 2 7
Thus A = 2/3 1 0 0 5/3 −11/3
1 6/5 1 0 0 −8/5
Writing UX = V, the given system becomes
1 0 0 𝑣 4
2/3 1 0 𝑣 = 5
1 6/5 1 𝑣 7
Solving this system, we have

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𝑣 =4 𝑣 = 7/3 𝑣 =1/5
Hence the original system becomes

3 2 7 𝑥 4
0 5/3 −11/3 𝑦 = 7/3
0 0 −8/5 𝑧 1/5
3x+2y+7z = 4: (5/3) x-(11/3)z = 7/3; (-8/5)z = 1/5
By back substitution, we have
Z = -(1/8) y = 9/8; and x = 7/8

Problems for discussion by the faculty in the class room

1. Ace Novelty wishes to produce three types of souvenirs: types A, B, and C. To manufacture a
type-A souvenir requires 1 minutes on machine I, 2 minute on machine II, and 3 minutes on
machine III. A type-B souvenir requires 2 minute on machine I, 3 minutes on machine II, and 4
minute on machine III. A type-C souvenir requires 3 minute on machine I and 4 minutes on
machines II and 1 minute on machine III. There are 14 min available on machine I, 20 min
available on machine II, and 14 min available on machine III for processing the order. How many
souvenirs of each type should Ace Novelty make in order to use all of the available time, write the
mathematical formulation to the problem and solve it by using LU decomposition method.

Home Work Problems

1. Ram, Raj, and Ravi go to a restaurant for lunch and order three different items. Ram orders
2 plates of fried rice, 3 plates of chicken pieces and 1-plate of curd rice. Raj orders 1 plate
of fried rice, 2 plates of chicken pieces and 3 plates of curd rice. Ravi orders 3 plates of
fried rice, 1 plate of chicken pieces and 2 plates of curd rice. Ram’s bill costs $9, Raj’s
costs $6, and Ravi’s costs $8. Determine plate cost of each item, write the mathematical
formulation to the problem and solve it by using factorization method.

2. Solve the following by using LU-decomposition method :

3 x  2 y  7 z  4, 2 x  3 y  z  5,3 x  4 y  z  7

Session 2: Taylor’s series for functions of two variables

Student outcome: To evaluate Taylor’s series expansion for functions of two variables

Let a function f(x,y) and all its derivatives up to nth order are finite and continuous at all points
(x,y) then, it can be written as an infinite power series in terms of (x-a) & (y-b) and is known as
Taylor’s series expansion of f(x,y) about the point (a,b) given by
f(x,y) = f(a,b)+ [(x-a)𝑓 (a,b)+(y-b)𝑓 (𝑎, 𝑏)]+
!
[(𝑥 − 𝑎) 𝑓 (𝑎, 𝑏) + 2(𝑥 − 𝑎)(𝑦 − 𝑏)𝑓 (𝑎, 𝑏) + (𝑦 − 𝑏) 𝑓 (𝑎, 𝑏)+---
!

Note : In the Taylor’s series when (a,b)→ (0,0) then the series is called Maclaurin’s series.
Maclurin’s series of f(x,y) about the origin is

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f(x,y)=f(0,0)+[x𝑓 (0,0)+y𝑓 (0,0)]+ [𝑥 𝑓 (0,0) + 2𝑥𝑦𝑓 (0,0) + 𝑦 𝑓 (0,0)]+---
!

Example : 1 Apply Taylor’s series to expand f(x,y)=x2+xy+y2 in powers of (x-1)and (y-2)

Sol: Differentiating f(x,y) = x2+xy+y2 partially w.r.to x & y we get
fx=2x+y; fyy=2; fxxx=0;
fy=x+2y; fyyy=0;
fxx=2; fyyx=0;
fxy=1; fxyy=0
The Taylor’s series expansion of f(x,y) in power of (x-1) and (y-2) is
f(x,y)=f(1,2)+[(x-1)𝑓 (1.2)+(y-2)𝑓 (1,2)]+ [(𝑥 − 1) 𝑓 (1,2) + 2(𝑥 − 1)(𝑦 − 2)𝑓 (1,2) +
!
(𝑦 − 2) 𝑓 (1,2)]+---
f(x,y)=7+4(x-1)+5(y-2)+ [2(𝑥 − 1) + 2(𝑥 − 1)(𝑦 − 2) + 2(𝑦 − 2) ] + 0+---
!

Problems for discussion by the faculty in the class room

1. Apply Taylor’s series to expand f(x,y)=x2-xy+y2 in powers of (x-1)and (y-2)

2. Expand the function f(x,y) =exlog(1+y) in terms of x and y up to the terms of second
degree using Taylor’s series
3. Expand f ( x, y )  sin x cos y in powers of x and y upto the terms of second degree
Home Work
1. Express the Taylor’s series expansion for f(x,y) = x2y+3y-2 in powers of (x-1) &(y+2)
2. Write the expansion of f ( x, y)  x 3  y 3  xy in powers of (x-1) and (y-2) by using
Taylor’s series

Session 3: Maxima and minima of functions of two variables

Student outcome: To determine the extreme values of function of two variables and Apply
Lagrange’s method to determine extreme values of function of three variables.

Definition : A function 𝑓(𝑥, 𝑦) is said to have a maximum or minimum at 𝑥 = 𝑎, 𝑦 = 𝑏,

according as 𝑓(𝑎 + ℎ, 𝑏 + 𝑘) < 𝑓(𝑎, 𝑏) (or) 𝑓(𝑎 + ℎ, 𝑏 + 𝑘) > 𝑓(𝑎, 𝑏).
Necessary conditions for extrema (maximum or minimum) of a function f (x,y)of two variables:
 The necessary conditions for f(x,y) to have maximum or minimum value at (a,b) are
f x (a, b)  0 & f y (a, b)  0
Stationary Value : 𝑓(𝑎, 𝑏) is said to be a stationary value of 𝑓(𝑥, 𝑦) if 𝑓 (𝑎, 𝑏) = 0 and
Working rule to find the maximum and minimum values of 𝒇(𝒙, 𝒚)
1. Find 𝑝 = and = 𝑞 and equate each to zero,
2. Solve these as simultaneous equations in 𝑥 and 𝑦. Let (𝑎, 𝑏), (𝑐, 𝑑) --- be the pairs of
values.
3. Calculate the values of 𝑟 = , 𝑠= , 𝑡= for each pair of values.
4.
i. If 𝑟𝑡 − 𝑠 > 0 and 𝑟 < 0 at (𝑎, 𝑏), 𝑓(𝑎, 𝑏) is maximum.
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 ii. If 𝑟𝑡 − 𝑠 > 0 and 𝑟 > 0 at (𝑎, 𝑏), 𝑓(𝑎, 𝑏) is minimum.
iii. If 𝑟𝑡 − 𝑠 < 0 at (𝑎, 𝑏), 𝑓(𝑎, 𝑏) is not an extreme value i.e., (𝑎, 𝑏) is a saddle
point.
iv. If 𝑟𝑡 − 𝑠 = 0 at (𝑎, 𝑏) then the case is doubtful and needs further
investigation.

Example 1. Determine the maximum and minimum of 𝑓(𝑥, 𝑦) = 𝑥 + 𝑦 + 6𝑥 + 12.

Sol. 𝑝 = 𝑓 = 2𝑥 + 6,𝑞 = 𝑓 = 2𝑦,𝑟 = 2,𝑠 = 0, 𝑡 = 2
For maximum and minimum
𝑓 = 0, 𝑓 = 0 gives 2𝑥 + 6 = 0 ---(1)
2𝑦 = 0 ---(2)
From (1) and (2)
𝑥 = −3 (or) 𝑦 = 0
The stationary point is (-3,0)
At (−3,0) ∶ (𝑟𝑡 − 𝑠 ) = 2 × 2 − 0 = 4 > 0 and 𝑟 = 2 > 0.
∴ at (-3,0), the function has minimum value.
Problems for discussion by the faculty in the class room

1. Determine the maximum and minimum of 𝑓(𝑥, 𝑦) = 𝑥 + 𝑦 + 8𝑥 + 4.

2. Divide 24 into two parts such that the continued product of the first and square of second
is maximum.
3. A rectangular box open at the top is to have volume of 32 cubic feet. Find the dimensions
of the box requiring least material for its construction.
Home Work
1. Determine the maximum and minimum of 𝑓(𝑥, 𝑦) = 𝑥 + 𝑦 + 6𝑥 + 12.
2. An electronics manufacture determines that the profit P (in dollars) obtained by producing
and selling x units of a DVD player and y units of a DVD recorder is approximated by the
model 𝑃(𝑥, 𝑦) = 8𝑥 + 10𝑦 − (0.001)(𝑥 + 𝑥𝑦 + 𝑦 ) − 10,000.Determine the
production level that produces maximum profit and also obtain the maximum profit
3. The profit obtained by producing x units of product A and y units of product B
approximately modeled by P ( x, y )  8 x  10 y  0.001( x 2  xy  y 2 ) .Estimate the
production level that produces a maximum profit.

Session 4:
Integral Calculus
Student outcome: Evaluate double integrals

Let R be a simple region defined by a  x  b and g1 ( x)  y  g 2 ( x) where g1(x) and g2(x)are

continuous functions on [a; b] and let f(x; y) be a function defined on R. Then
b  g2  x  
I R    f  x, y  dy  dx

 
a  g1  x  
is called an iterated integral of f over R.

6
Similarly, if R is a simple region defined by c  y  d and h1 (y)  x  h2 (y) , where h1(y) and
h2(y) are continuous functions on [c; d] and let f(x; y) be a function defined on R. Then
d  h2  y  
I R    f  x, y  dx  dy

 
c  h1  y  
f
is called an iterated integral of over R.

4 4 x
Example 1: Evaluate the integral    xy  dx dy
1 0

Solution:
Given that the integral
4  4 x 4 x
4 4 x  4
 y2  4
4  x
1 0  
xy dx dy  1  0   
 xy dy  dx  1  2 
x dx  1 x  2  dx
  0
4
 x2 
4
 x3  32 5 9
   2 x   dx   x 2      .
1
2   6 1 6 6 2
x
Example(2). Compute the integral  dA where D is the triangle with vertices
D y  1
 0, 0  , 1, 1 , and  2, 0  .
Solution:

7
The integral is evaluated as follows:
1 2 y 2 y
1  x2 
1 1
x x 1 1 
 dA    dxdy     dy    2  y   y 2  dy
2

D
y 1 0 y
y 1 0
y 1  2  y 2 0 y 1  
1 1
1 1 1 y
   4  4 y  dy  2 dy
2 0 y 1 0
1 y
1
 2 
 1  dy  2  2 log 1  y   y  0  4 log 2  2.
1
 2 
0
1 y 

Problems for discussion by the faculty in the class room

3 2
1. Evaluate the integral   x  x  y  dydx
x  0 y 1

2. Use double integration to find ∬ 𝑥𝑦 𝑑𝑥𝑑𝑦 over the positive quadrant of the circle x2+y2=a2
3. Sketch the region R in the xy-plane bounded by the curves y 2  x and y  x find its area.

Home Work
1. Determine the area of the region bounded by the Parabolas y 2 = x and x2 = y.
2. Evaluate∫ ∫ 𝑒 𝑑𝑦𝑑𝑥
3. Evaluate the area of the region bounded by x-axis, x=2a and the curve x 2=4ay.

Session-5:
Student outcome: Evaluate double integrals by the change of order of integration

Here x and y are independent variables, one could just as well exchange the order of integration,
deriving first the area of a strip parallel to the x-axis, then summing these along y, as shown in the
figure.

1. Change the order of integration and evaluate the iterated integral:

8
 1 1

   xy  sin  x   dxdy
4

0 y1/3

Sol: From the figure we see that the region D is bounded above by y  x3 and below by y  0 .
The projection of D onto the x-axis the interval 0  x  1. Using the order of integration dydx we
have

5 1
  cos 1
16 4

Evaluate double integral by transforming into polar coordinates

Change two dimensional Cartesian to polar coordinates. The Jacobian is r which is always positive
so that the modulus is r. Thus the element of area is dA  rdrd . This is readily seen from the
diagram.

9
Problems for discussion by the faculty in the class room

1. Change the order of integration and hence evaluate ∫ ∫ (𝑥 + 𝑦 )𝑑𝑦𝑑𝑥

a a2  y2

  x  y 2  dydx by changing to polar coordinates.

2
2. Evaluate
0 0

3. Evaluate ∫ ∫ ( )
𝑑𝑥𝑑𝑦 by changing to polar coordinates.

Home work
1 1 x 2
1. Apply change the order of integration and evaluate  
0 0
y 2 dx dy.

4a 2 ax
2. Evaluate the integral by reversing the order of integration   dy dx.
0 x 2 /4 a

3. Compute the integral  xy dxdy over

R
the quarter circle [use Polar or Cartesian]

R   x, y  : 0  x, 0  y , x 2  y 2  1
1 y
4. Evaluate the iterated integral by converting to polar coordinates 
0 y
x 2  y 2 dxdy over

the region R bounded by the circle x 2  y 2  9 in the first quadrant.

Session-6
Student outcome: Evaluate triple integrals

Consider a function f  x, y, z  at every point of the three dimensional finite region V. Divide V
into n elementary volumes  V1 ,  V2 ,..........,  Vn . Let  xr , yr , zr  be any point within the rth sub-
division  Vr . Consider the sum
n

 f  x , y , z V
r 1
r r r r

The limit of the sum , if it exits , as n   and  Vr  0 is called the triple integral of f  x, y, z 
over the region V and is denoted by
 f  x, y, z  dV
V

For purposes of evaluation, it can also be expressed as the repeated integral

x2 y2 z2
   f  x, y , z  dzdydx
x1 y1 z1

The order of integration may be different for different types of limits.

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Example (1): Find the volume of the solid in the first octant bounded by the planes
x  0, y  0, z  0, x  2 y  z  6
Solution: The region is plotted below.

Therefore the volume V   dV

D
The region D bounded below the plane z  0 and above by the plane z  6 - x - 2 y.
The projection of D onto the xy-plane is the triangular region
 6 x 
 x, y  : 0  y  , 0  x  6 .
 2 
Using the order of integration dz dy dx, the volume V is:
6  6 x / 2 6 x  2 y
V   dV     dz dy dx
0 0 0
D
 6 x  / 2
 z 0
6 6 x  2 y
  dy dx
0 0
 6 x / 2
 6  x  2 y  dy dx
6
 
0 0
6  6 x / 2
   6 y  xy  y 2  dx
0 0

6   6  x   6  x   6  x 2 
 6    x    dx
  2   2   2  
0

6 x 
2
3
    x 2  9 x  dx
0
 4 2 
6
1 3 
  x 3  x 2  9 x   18.
12 2 0
Evaluate triple integrals by changing into cylindrical and spherical coordinates

Cylindrical Polar Coordinates:

Let us choose x  r cos  ; y  r sin  ; z  z , and then the differential volume is
  x, y , z 
dV  dr d dz  r dr d dz
 r ,  , z 

11
Spherical Polar Coordinates:
Let us choose x  r sin  cos  ; y  r sin  sin  ; z  r cos  ,then the differential volume is
  x, y , z 
dV  dr d d  r 2 sin  dr d d .
 r ,  ,  

Problems for discussion by the faculty in the class room

1 1 1
1. Evaluate  x 0 
y 0 
z 0
x 2 yz dx dy dz
1 x x y
2. Evaluate    ( x  y  z ) dz dy dx
x  0 y 0 z  0

3. Determine the volume of the region bounded by the plane x + y + z = 8 that lies in the
first octant.
4. Evaluate the volume of the tetrahedron bounded by the planes x=0,y=0,z=0,x+y+z=1.

Home Work
a x x y

1. Evaluate   
0 0 1
e x y  z dx dy dz
y
9 3 y 2 9 x 2

2. Evaluate Obtain the triple integral   

y 0 x 0 z 0
zdzdxdy

3. Evaluate the following integral.  8xyz dV , where B : 2  x  3,1  y  2 and 0  z 1

B

12
Vector Calculus
Session No: 7
Student Outcome: A student can able to apply the knowledge of gradient , divergent and curl to
solve real time applications in 2D plane and 3D space.
VECTOR DIFFERENTIATION:
Vector differential operator () :The vector differential operator  (read as del) is defined as
  
i  j k
x y z

We now define some quantities known as gradient, divergence and curl.

Gradient of a scalar point function :

Let f ( x, y, z ) be a scalar point function defined in some region of space. Then the vector
f f f
function i  j k is known as the gradient of f and is denoted by grad f or f
x y z
   f f f
f = (i  j k )f  i  j k
x y z x y z
Example-1. If f  x, y, z   3 x y  y z , find grad f at the point 1,2,1
2 3 2

  
Solution: we know f = (i  j  k ) f
x y z
  
= (i  j  k )3 x 2 y  y 3 z 2 
x y z
 
= 6 xi  3 x  3 y z j  2 y 3 zk
2 2 2

Putting x  1, y  2, z  1 , weget f =  12i  9 j  16k

Directional derivative:
Directional derivative of scalar point function: Let f  x, y, z  define a scalar field in a region
R and P be any point in the region. The directional derivative of a scalar field f at a P  x, y , z  in
the direction of a unit vector â is given by f  aˆ .
NOTE: Therefore the directional derivative is maximum along the normal to the surface. Its
maximum
value is f .
Example 2. Find the directional derivative of f  xy  yz  zx in the direction of vector
i  2 j  k at the point (1,2,0).
Sol: Given f  xy  yz  zx .
  
f = (i  j k )f
x y z
= ( y  z )i  ( x  z ) j  ( x  y ) k
If a is the unit vector in the direction of the vector i  2 j  k , then

13
 i  2 j  2k i  2 j  2k
aˆ  
1 4  4 3

There for Directional derivative of f along the given direction =  f .a
i  2 j  2k
 . ( y  z )i  ( x  z ) j  ( x  y ) k at (1,2,0)
3
= 10/3.
Example-3 Evaluate the angle between the normal’s to the surface xy  z 2 at the points
(4,1,2) and (3,3,3) .

Solution Given surface is f ( x, y, z )  xy  z 2

Let n1 and n 2 be the normal’s to the surface at (4,1,2) and (3,3,3) respectively
Differentiating partially, we get
f f f
y;  x ;   2z
x y z
Grad f = yi+xj - 2zk
n1 ( gradf ) at ( 4,1,2)  i  4 j  4k

n1 ( gradf ) at (3,3,3)  3i  3 j  6k
Let  be the angle between the two normals
n1 .n 2 (i  4 j  4k ).(3i  3 j  6k ) 9
cos    =
n1 n 2 1  16  16 9  9  36 33 54

Divergence of a vector point function:

Let f be any continuously differentiable vector point function. Then the divergence of f ,is
written as div f or . f

i.e div f = 𝑖 + 𝑗 + 𝑘 ̅ (𝑖 . 𝜕𝑓 ̅/𝜕𝑥 + 𝑗 . 𝜕𝑓 ̅/𝜕𝑦 + 𝑘. 𝜕𝑓 ̅/𝜕𝑧)

. 𝑓=
 f f f
If f = f1i  f 2 j  f 3 k then divf  1  2  3 .
x y z
Solenoid vector:
A vector point function f is said to be solenoidal if div f =0

Example -4 If f  xy i  2 x yzj  3 yz k find div f at (1,1,1)

2 2 2

Given f  xy i  2 x yzj  3 yz k . Then

2 2 2
Solution

14
 f f f
div f     y 2  2 x 2 z  6 yz
x y z
(div f ) at (1,1,1)  9

Curl of a vector:
Def: Let f be any continuously differentiable vector point function. Then the vector function
f f f
defined by i   j k is called curl of f and is denoted by curl f or   f
x y z
i j k
f f f f   
Curl f = i   j k = i  =
x y z x x y z
f1 f2 f3
Irrotational vector:
A vector is said to be Irrotational if Curl f =0
If f is irrotational , there will always exits a scalar function  ( x, y, z ) such that f = grad  .
This  ( x, y, z ) is called scalar potential of f .
It is easy to prove that, f = grad  , then Curl f =0.

Example -5 If f  xy 2i  2 x 2 yzj  3 yz 2 k Find Curl f at the point (1,-1.1)

Solution Let f  xy 2i  2 x 2 yzj  3 yz 2 k . Then
i j k
  
Curl f =   f =
x y z
xy 2 2 x yz  3 yz 2
2

     
= i[ ( 3 yz 2 )  ( 2 x 2 yz )] - j[ (3 yz 2 )  ( y 2 x)] + k [ ( 3 yz 2 )  ( 2 x 2 yz ) ]
y z sx z y z
=  (3 z 2  2 x 2y )i  (4 xyz  2 xy )k

 curl f at (1,1,1)   i  2k .

Example -6 Prove that Curl r  o

Solution Let r  xi  yj  zk

Curl curl r   i  (r )   (i  i)  o
x

Example -7 Prove that 𝑑𝑖𝑣 𝑟̅ = 3

15
Solution Let r  xi  yj  zk
𝑑𝑖𝑣 𝑟̅ = + + =3

Problems for discussion by the faculty in the class room

1. Determine a vector normal to the surface f  x 2  2 y 2  4 z 2 at (1, 1, -1)
2. The temperature at the point  x, y, z  in space is given by T  x, y, z   x  y  z. A
2 2

mosquito located at  4, 4, 2 desires to fly in such a direction that it gets cooled faster.
Determine the direction in which it should fly.
3. Obtain the directional derivative of   xy  yz  zx at A in the direction of AB where A=
(1,2,-1) , B=(1,2,3) .
4. Find div f and Curl f where f  grad ( x 3  y 3  z 3  3 xyz ) .
Home Work
1. Find a vector normal to the surface x 2  y 2  2 z 2  26 at the points (2,2,3).
2. If T is the temperature field given byT = x − y + xz , compute the gradient of
temperature T at the point 1,1,2  .
3. Evaluate the angle between the normal’s to the surface xy  z 2 at the points
(4,1,2) and (3,3,3) .
4.
5. Find the constants a, b and c if the vector
f  (2 x  3 y  az ) i  (bx  2 y  3 z ) j  (2 x  cy  3 z ) k is irrotational.

6. Determine p, if F   x  3 y  i   y  2 z  j   x  pz  k is solenoid.
7. Obtain the curl and divergent of the vector field V   x 2 y 2  z 3  i  2 xyz j  e xyz k

Session No: 8
Student Outcome: A student can able to apply the knowledge of Line Integral to find out the
total workdone by a particle.

VECTOR INTEGRATION
Integration of vectors
( )
If two vector functions F (t) and G (t) be such that = F(t), then G (t) is called an integral of F
(t) with respect to the scalar variable t and we write
𝐹(𝑡)𝑑𝑡 = 𝐺(𝑡) + 𝑐
Where c is an arbitrary constant vector
This is called the indefinite integral of F (t) and it definite integral is

16
 𝐹(𝑡) = [𝐺(𝑡) + 𝑐] = 𝐺(𝑏) − 𝐺(𝑎)

Ex: -Given𝑓(𝑡) = (5𝑡 −)𝑖 + 6𝑡 𝑗 − 7𝑡𝑘 , 𝑒𝑣𝑎𝑙𝑢𝑎𝑡𝑒 ∫ 𝑓(𝑡)𝑑𝑡

Sol:-∫ 𝑓(𝑡)𝑑𝑡 = ∫ [(5𝑡 − 3𝑡)𝑖 + 6𝑡 𝑗 − 7𝑡𝑘)] 𝑑𝑡
5𝑡 3𝑡 3 7𝑡
= − 𝑖+ 𝑡 𝑗− 𝑘
3 2 2 2
5 3 3 7
= (64 − 8) − (16 − 4) 𝑖 + (256 − 16)𝑗 − (16 − 4)𝑘
3 2 2 2
226
= 𝑖 + 360 𝑗 − 42 𝑘
3
LINE INTEGRAL
𝐼𝑓 𝑎 𝑣𝑒𝑐𝑡𝑜𝑟 𝑓𝑢𝑛𝑐𝑡𝑖𝑜𝑛 𝐹(𝑅) = 𝑓𝑖 + 𝑔𝑗 + ℎ𝑘
𝑊ℎ𝑒𝑟𝑒𝑓, 𝑔 , ℎ𝑎𝑣𝑒 𝑓𝑢𝑛𝑐𝑡𝑖𝑜𝑛𝑠 𝑜𝑓 𝑥 , 𝑦, 𝑧 𝑎𝑛𝑑 𝑑𝑅 = 𝑑𝑥𝑖 + 𝑑𝑦𝑗 + 𝑑𝑧𝑘, Then
∫ 𝑑𝑅 = ∫ (𝑓𝑑𝑥 + 𝑔𝑑𝑦 + ℎ𝑑𝑧)
is called the line integral (or tangential line integral) of F over c, where c is any curve in space
Work done by a Force
If F represents the force acting on a particle moving along an arc AB, then the total work done by
F is given by the line integral ∫ 𝐹 . 𝑑𝑅
𝐼𝑓 𝑓 ̅ = (5𝑥𝑦 − 6𝑥 )𝑖 + (2𝑦 − 4𝑥)𝑗,
Evaluate  F .dR along the curve C in the xy – plane, y = x 3 from (1, 1) to (2, 8).
C

Sol:- 𝑓 = (5𝑥𝑦 − 6𝑥 )𝑖 + (2𝑦 − 4𝑥)𝑗

𝑑𝑅 = 𝑑𝑥𝑖 + 𝑑𝑦𝑗
∴ 𝐹 . 𝑑𝑅 = (5𝑥𝑦 − 6𝑥 )𝑑𝑥 + (2𝑦 − 4𝑥)𝑑𝑦
Equation of the cure is y=𝑥 ⟹ 𝑑𝑦 = 3𝑥 𝑑𝑥
∴ 𝑜𝑛𝑡ℎ𝑒𝑐𝑢𝑟𝑒,
𝐹 . 𝑑𝑅 = 5𝑥. 𝑥 − 6𝑥 𝑑𝑥 + (2. 𝑥 − 4𝑥) . 3𝑥 𝑑𝑥
= (5𝑥 − 6𝑥 + 6𝑥 − 12𝑥 )𝑑𝑥

∴ ∫ 𝐹 . 𝑑𝑅 = (5𝑥 − 6𝑥 + 6𝑥 − 12𝑥 ) 𝑑𝑥

[𝑥 ] − [2𝑥 ] , [𝑥 ] − [3𝑥 ]
= (32 − 1) − (14) + (64 − 1) − (45)
= 35
Problems for discussion by the faculty in the class room

1. Obtain the work done in moving a particle in the force field F = (3x+6y)I – 14yz J+ 20 xz 2
K , along (i) the straight line from 0,0,0 to (1,1,1) (ii)the curve defined by
x 2  4 y,3 x 3  8 z from x  0 to x =1

17
 2. Using the line integral, compute the work done by the force F = 3xy I-5z J +10x K when it
moves a particle from the point 0,0,0 to the point  2,1,1 along the curve x = t, y = t2 +1,
z = t3 from t = 0 to t = 1
Home Work

 ( xy  z
2
1. Evaluate )dz where C is the arc of the helix x  cos t , y  sin t , z  t which joins
C

the points 1,0,0 and  1,0,   .

  
2. Evaluate  F .d r where F  xy i   4 x  2 y  j and C is the line segment from
C

 4, 3 to  7,0 .
  
3. Evaluate  F .d r where F   x 3  y  i   x 2  7 x  j and C is the portion of y  x 3  2
C

from to .

SESSION - 9
GREEN’S THEOREM IN THE PLANE
(Transformation between line integral and double integral)
If R is a closed region in xy plane bounded by a simple closed curve C and if M and N are
continuous functions of x and y having continuous derivatives in R then

∂N ∂M
∫ (M dx + N dy) = ∫ ∫ − dxdy
∂x ∂y
Where C is traversed in the positive (antilock wise) direction.
Problems for discussion by the faculty in the class room
1) Apply Green’s theorem to evaluate the integral ∫ [(xy + y )dx + x dy] where C is
bounded by y=x and y=x2.

2) Evaluate ∫ (x + xy)dx + (x +y )dy

Where C is the square formed by the lines x = y = ±1 and x = ±1

Home Work

1. Apply Green’s theorem to evaluate the integral ∫ (2x − y )dx + (x + y ) dy where

C is the boundary of the area enclosed by the x-axis and upper half of the circle x + y =
a .

18
 2. Apply Green’s theorem, evaluate  [( y  sin x)dx  cos x dy] where C is the plane triangle
C

2
Enclosed by the lines y =0, x = π/2 and y  x

3. A Vector Field is given by F  Sin y i  x(1  Cos y ) j Evaluate the line integral over the
.

circular path x 2  y 2  a 2 , z  0. by using Green’s Theorem.

Session 10 SURFACE INTEGRALS

Consider a continuous function F (R) and a surface S. If N is a Unit normal (outward) to the surface
at any point, then ∫ 𝐹. 𝑁𝑑S is the normal surface integral (or surface integral ) of F(R) over S.

(1) 𝐼𝑓 𝐹 = 2𝑦𝑖 − 3𝑗 + 𝑥 𝑘 𝑎𝑛𝑑 𝑆 𝑖𝑠 𝑡ℎ𝑒 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑜𝑓 𝑡ℎ𝑒 𝑝𝑎𝑟𝑎𝑏𝑜𝑙𝑖𝑐 𝑐𝑦𝑙𝑖𝑛𝑑𝑒𝑟

𝑦 = 8𝑥 𝑖𝑛 𝑡ℎ𝑒 𝑓𝑖𝑟𝑠𝑡 𝑜𝑐𝑡𝑎𝑛𝑡 𝑏𝑜𝑢𝑛𝑑𝑒𝑑 𝑏𝑦 𝑡ℎ𝑒 𝑝𝑙𝑎𝑛𝑒𝑠 𝑦 = 4 𝑎𝑛𝑑 𝑧 = 6, 𝑠ℎ𝑜𝑤 𝑡ℎ𝑎𝑡

∫ 𝐹. 𝑁𝑑𝑠 = 132

Sol:-
A unit normal to the surface 𝑦 = 8𝑥 𝑖𝑠
−8𝑖 + 2𝑦𝑗 −4𝑖 + 𝑦𝑗
𝑁= =
4𝑦 + 64 𝑦 + 16
−8𝑦 − 3𝑦 −11𝑦
∴ 𝐹. 𝑁 = =
𝑦 + 16 𝑦 + 16

𝑁. 𝑖 =

𝐶𝑜𝑛𝑠𝑖𝑑𝑒𝑟 𝑡ℎ𝑒 𝑝𝑟𝑜𝑗𝑒𝑐𝑡𝑖𝑜𝑛 𝑜𝑓 𝑡ℎ𝑒 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑜𝑛 𝑦𝑧 𝑝𝑙𝑎𝑛𝑒

∴ 𝑑𝑠 = | .|
.= =

−11𝑦
∴ 𝑅𝑒𝑞𝑢𝑖𝑟𝑒𝑑 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑖𝑛𝑡𝑒𝑟𝑔𝑟𝑎𝑙 = 𝑑𝑦𝑑𝑧
4

11 𝑦
= − [𝑧]
4 2

−11 16
= . .6 = −132
4 2
Absolute value = 132

19
(2) Evaluate ∫ 𝐹. 𝑁𝑑s where 6𝑧𝑖 − 4𝑗 + 𝑦𝑘 𝑎𝑛𝑑 𝑆 𝑖𝑠 𝑡ℎ𝑒 𝑝𝑜𝑟𝑡𝑖𝑜𝑛 𝑜𝑓 𝑡ℎ𝑒 𝑝𝑙𝑎𝑛𝑒 2𝑥 +
3𝑦 + 6𝑧 = 12 𝑖𝑛 𝑡ℎ𝑒 𝑓𝑖𝑟𝑠𝑡 𝑜𝑐𝑡𝑎𝑛𝑡

Sol:-
Let ∅ = 2𝑥 + 3𝑦 + 6𝑧 − 12
∴ 𝐴 𝑛𝑜𝑟𝑚𝑎𝑙 𝑡𝑜 𝑡ℎ𝑒 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑖𝑠 ∇∅ = 2𝑖 + 3𝑗 + 6𝑘
∇∅ 2𝑖 + 3𝑗 + 6𝑘 2𝑖 + 3𝑗 + 6𝑘
∴ 𝐴 𝑢𝑛𝑖𝑡 𝑛𝑜𝑟𝑚𝑎𝑙 𝑡𝑜 𝑡ℎ𝑒 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 𝑖𝑠 𝑁 = = =
|∇∅| √4 + 9 + 36 7
( )
∴ 𝐹. 𝑁 = =

Consider the projection of the surface on yZ plane

𝑑𝑦𝑑𝑧 7
∴ 𝑑𝑠 = = 𝑑𝑦𝑑𝑧
|𝑁. 𝑖| 2
2 7
∴ 𝐹. 𝑁 𝑑𝑠 = (6𝑧 − 6 + 3𝑦). 𝑑𝑦𝑑𝑧
7 2
= (6𝑧 + 3𝑦 − 6)𝑑𝑦𝑑𝑧
𝑇ℎ𝑒 𝑒𝑞𝑢𝑎𝑡𝑖𝑜𝑛 𝑜𝑓 𝑡ℎ𝑒 𝑠𝑢𝑟𝑓𝑎𝑐𝑒 2𝑥 + 3𝑦 + 6𝑧 = 12 𝑜𝑛 𝑡ℎ𝑒 𝑦𝑧 𝑝𝑙𝑎𝑛𝑒 𝑖𝑠 3𝑦 + 6𝑧 = 12
∴ y ∶ 0 → 4 and
12 − 3𝑦
𝑧 ∶ 0→
6

∴ ∫ 𝐹. 𝑁𝑑𝑠 =∫ ∫ (6𝑧 + 3𝑦 − 6)𝑑𝑧𝑑𝑦

= [3𝑧 + 3𝑦𝑧 − 6𝑧] 𝑑𝑦

12 − 3𝑦 (12 − 3𝑦) (12 − 3𝑦)

= 3𝑦 −6 𝑑𝑦
6 6 6

=∫ 3𝑦 − 𝑦 𝑑𝑦 = − = . 16 − =8

(3) Evaluate  F  dS
S
if F  yz i  2 y 2 j  xz 2 k and S is the surface of the cylinder

x 2  y 2  9 contained in the first octant between the planes z  0 and z  2 .

Sol.: Given that F  yz i  2 y 2 j  xz 2 k

20
   
Let   x 2  y 2  9  0 .Then  2x ,  2y , 0
x y z


 grad    i  2 xi  yj 
x

 grad   2 x 2  y 2  23  6

grad  xi  yj
 N  Unit normal  
grad  3

Let R be the projection of S on yz-plane

dydz xyz  2 y 3 dydz

Then  F  N dS   F  dS   F  N R
= 
N i R  
x
3
For the surface in the yz-plane x=0  y=3,z  0 and 2.

3  2y3 
2
  F  N dS  3   yz  dydz
0
0  x 

2
2y3  3
y3
= 3   yz 
3
dydz  6 y 3  12
0 
0 9  y 
2  0 0 9  y 2 dy (1)

3
y3
Putting y  3 Sin  , we get 
0 9 y 2
dy  18 (2)

  F  N dS  18  12  18  234, using (1) and (2).

Problems for discussion by the faculty in the class room

     
1) Evaluate  A.Nds where A  zi  xj  3 y 2 zk and S is the surface of the cylinder
S

x 2  y 2  16, included in the Ist octant between Z = 0 and Z = 5

3
2) Show that  F .Nds  2
S
where F  4 xzi  y 2 j  yzk and `S’ is the surface of the cube

bounded by the planes x=0,x=1,y=0,y=1 and z=0,z=1.

Home Work

1. Evaluate  F .Nds ,
S
where F  2 yxi  yzj  x 2 k , over the surface, S of the cube

bounded by the coordinate planes and planes x=a, y=a and z=a.

21
 2. Evaluate the flux of the vector field A   x  2 z  i   x  3 y  z  j   5 x  y  k
through the upper side of the triangle ABC with vertices at the points A
(1,0,0),B(0,1,0) and C(0,0,1).
3. Evaluate  A.Nds where
S
A  18 zi  12 j  3 yk and S is that part of the plane

2x+3y+6z=12 which is located in the first octant.

4. If F  yz i  zx j  xy k , evaluate   F  N dS over the surface x 2  y 2  z 2  1 in
the first octant.

SESSION-11
STOKE’S THEOREM
If S be a open surface bounded by a closed non intersecting curve C and F is any differentiable
vector point function then ∫ F. dr = ∫ curl F . n ds
Where n is a unit outward normal at any point on the surface.

Problems for discussion by the faculty in the class room

  
1. Apply Stoke’s theorem to evaluate ∫ F. dr where F  ( x 2  y 2 )i  2 xyj where C is the
boundary of the rectangle bounded by the lines x = ±a, y = 0, y = b.

2. Using Stoke’s theorem evaluate  [( x  y )dx  (2 x  z )dy  ( y  z )dz ] where C is the boundary
C

of the triangle with vertices (2,0,0), (0,3,0) and (0,0,6)

Home Work
1. Apply Stokes theorem for the function 𝐹 = 𝑥 𝚤̅ + 𝑥𝑦 𝚥⃐ integrated round the square in the
plane z=0 whose sides are along the lines x=0, y=0, x=a, y=a.
2. Using Stoke’s theorem evaluate the integral ∫ F. dr where
𝐹 = 2𝑦 𝚤̅ + 3𝑥 𝚥̅ − (2𝑥 + 𝑧)𝑘 where C is the boundary of the triangle whose vertices
are (0, 0, 0),(2, 0, 0),(2, 2, 0).

22
Ordinary Differential Equations and its applications
 Newton’s law of cooling

 The laws of natural growth and decay

Differential Equations of First Order
LEARNING OBJECTIVES

After reading this chapter, the student will be able to understand:

 Formation of differential equation

 Methods of solving first order differential equations in different standard forms

 Modeling and solving the physical problems through first order differential
equations

Introduction

The mathematical formulation of problems in engineering and science usually leads to equations
involving derivatives of one or more unknown function.Such equations are called differential
equations.
In this chapter we shall consider various physical and geometrical problems that lead to differential
equations, with emphasis on modelling, i.e., the transition from the physical situation to a
mathematical model. Here we explain most important analytical methods for solving such
equations.
The study of a differential equation consists of the following three phases:
(i) Formulation of differential equation from the given physical situation, called modelling.

(ii) Solving of this differential equation, evaluating the arbitrary constants from the given
initial or boundary conditions.

(iii) Physical interpretation of the solution.

Some of the common applications where the differential equations would be used:
 Modeling biological growth, radioactivity and carbon dating

 Oscillations of mechanical and electrical systems.

 Used to understand and predict the spread of diseases.

 To predict weather changes.

 To detect early signs of heart disease.

23
  Conduction of heat.

 Chemical reactions and mixture problems.

 Geometrical problems.

 Problems in economics.

Models in first order differential equations

Consider the motion of the body of mass m along a straight line, which we designate as an x-
axis .Let the mass be subjected to a force F(t) along that axis, where t is time. Then according to
the Newton’s second law of motion
d 2x
m 2  F (t ). (1)
dt
where x(t) is the displacement of mass measured from the origin.
If we know the applied force F(t) and wish to determine the displacement x(t) , then it is needed
to solve the differential equation.
Suppose that F(t)= F0 is constant ,then integrating (1) with respect to t on both sides we get

F0 2
mx= t  At  B , where A and B are orbitary constants.
2
But most differential equations cannot be solved easily, i.e. by integration.
Suppose the mass is connected to a coil spring that supplies a restoring force proportional to the
displacement with constant of proportionality k, and then the differential equation (1) becomes
d 2x
m 2   kx  F (t )
dt
Integrating this
dx
m   k  x (t ) dt   F (t ) dt  A ,
dt
where A is the constant of integration. Since F(t) is a prescribed function ,the integral of F(t) can
be evaluated. But since x(t) is the unknown ,the integral of x(t) cannot be evaluated and we cannot
proceed with our solution by integration.

Session 12

The laws of natural growth and decay

24
In many natural phenomena, quantities grow or decay at a rate proportional to their size/amount.
If x(t) is the value of some quantity y at time ‘ and if the rate of change of x with respect to ‘ ’
is proportional to x , then
dx
 kx (where is a constant)
dt
This equation is called the law of natural growth (if ) or the law of natural decay (if
).

Note: If x(t) denotes the amount of a radioactive substance present at time t, then

dx
 kx (k  0)
dt

he rate dx/ dt is negative, since x is decreasing. The positive constant k is called the rate constant
for the particular radioisotope. The solution of this separable first-order equation is

x  x0 e  kt (1)

where x0 denotes the amount of substance present at time t = 0. The graph of this equation (Figure
3.1)is known as the exponential decay curve:

The relationship between the half-life (denoted by T1/2) and the rate constant k can easily be found.
Since, by definition, x = ½ x0 at t = T1/2, equation (1) becomes

1  kT ln(2)
x0  x0 e 1 2   kT1 2  ln(1 / 2)  k 
2 T1 / 2

25
It shows that the half-life time and rate constant are inversely proportional, the shorter the half-
life, the greater the rate constant, and consequently, the more rapid in decay.

Example 1 It is known that at any instant of time, radium decays at a rate proportional to the
amount present. If the initial mass of radium is m0 find an expression for the mass of radium
at time t.
Solution Let be mass of radium at time ‘ ’. Then the rate of change of its mass is

= , (1)

where is the constant of proportionality and is the mass. The -ve sign is taken since the
mass is decreasing with time.
Solving equation (1) by separating variables , we get
e ⇒ m(t)=Ae-kt (2)

But it is given that at time 0. Hence 0

Therefore, equation (2) becomes m(t )  m 0 e  kt

This is the expression for the mass of radium remain at time .

Example 2 It is given that the rate of decay of radium varies as its mass at that time. Assuming
that the half life of radium is 1600 years, find what percent of the mass 0 will remain after
200 years.
Solution Let be the mass (in gms) of radium at any time ( in years)

dm
Then the differential equation representing the decay is  km (1)
dt
Separating the variables and integrating, we obtain

But at time 0. Hence c   log( M 0 )

(2)

Since the half life of radium is 1600 years

 M 2 1
-1600k= log 0   log(1 / 2)  k  ln(2)
 M 0  1600

Substituting of value of k in equation (2)

26
Let the mass of radium after 200 years be 1 gms. Then

 200 M  M 1
ln(2)  Log  1   1  8
1600  M0  M0 2
The percentage of mass that remains after 200 years is given by
M1
100  (100) 2 1 8 =91.7 %
M0

Example 3 The number of bacteria in a culture grew at a rate proportional to . The value
of was initially 100 and increased to 332 in one hour. What was the value of after 1 ½
hours?

Solution It is given that

dN dN
 kN   kdt
dt N

On integrating, logN=kt+c  N  Aekt (1)

Given that, when

(2)

Also, given that, when

332
k
 k  log
100
t
 332 
Substituting the value of k in (2) , N  100 
 100 
3
 332  2
When t=3/2 ,the number of bacteria is N  100   604.9  605
 100 
Example 4 A fragment of bone is discovered to contain 20% of the usual C-14 concentration.
Estimate the age of the bone, given that half life time of C-14 is 5730 years.
Solution The relative amount of14C in the bone has decreased to 20% of its original value (that
is, the value when the animal was alive). Thus, the problem is to calculate the value of t at which
x( t) = 0.20 xo (where x = the amount of14C present). Since

27
 ln 2 ln 2
k   the exponential decay equation 
T1 5730
2

  ln 2/5730  t ln 2
0.20 x0  x0 e  ln(0.20)   t
5730
 ln(0.20)
 t  5730  t  13, 300 years
ln 2
Example 5 When a chicken is removed from an oven, its temperature is measured at 300 0F.
Three minutes later its temperature is 200o F. How long will it take for the chicken to cool off to
a room temperature of 70oF.
Solution Given that Ts  70, T0  300
By Newton’s law of cooling T  Ts  (T0  Ts )e  kt
Substituting the values of Ts and T0
T  70  (230)e  kt
Given that T=200 when t=3
13 1  23 
 200  70  230e 3k  e 3k   k  ln   0.19018
23 3  13 
T  70  (270)e 0.19018t
Now, T  70 when T  ∞
Thus the chicken cools off to room tempeture after a long period of a time.
Example 6 The population of a community is known to increase at a rate proportional to the
number of people present at a time t. If the population has doubled in 6 years, how long it will
take to triple?
Solution Let N(t) denote the population at time t. Let N 0 denote the initial population
(population at t=0).
dN
 kN(t )
dt
Solution is N(t)= N 0 ekt (1)
Since the population has doubled in 6 years

2 N 0 = N 0 e6k  k  (log 2) / 6
We have to find the volume of t, when N  3N 0
Substituting N  3N 0 in (1) we get
3 N 0 = N 0 ekt

1 ln 3
t  ln 3  6  9.6 years
k ln 2

Example7 Let population of a country be decreasing at the rate proportional to its population. If
the population has decreased to 25% in 10 years, how long will it take to become half?

28
 dN
Solution This phenomenon can be modeled by  kN (t )
dt
Its solution is given by
N(t)= N 0 ekt, where N 0 is the initial population (1)
Given that for t=10, N(10)= N 0 / 4
Substituting this in (1),
1 1
N 0 4 = N 0 e10k  e10k= k= ln (1/4)
4 10
Substituting the value of k in (1)
1
ln( t / 4 )
N (t )  N 0 e 10
The time require for the population to become half is given by
1
ln( t / 4 ) ln(1 / 2)
N 0 / 2  N 0 e 10 t   24 years ( approx ).
(1 / 10) ln(1 / 4)
Example8 A radioactive isotope has an initial mass 200mg, which two years later is 150 mg.
Find the expression for the amount of the isotope remaining at any time. What is its half-life?
Solution Let m be the mass of the isotope remaining after t years, and let -k be the constant of
dm
proportionality. Then the rate of decomposition is modeled by = - km,
dt
where minus sign indicates that the mass is decreasing.
Solving this equation m = c e-kt ,
To find c , recall that m =200 when t=0. Putting these values of m and t we get
200 = c e-k.0 = c.1 or c=200
 m = 200e-kt (1)
The value of k may now be determined by substituting t=2, m=150 in (1)
150 = 200 e-2k or e
2k  3 or –2k=ln 3
4 4
1 4
This gives, k  ln = 0.1438  0.14
2 3
The mass of the isotope remaining after t years is then given by m(t) =200e -0.14t.
The half-life t1/2 is the time corresponding to m=100mg is given by
1 1 0.693
100 = 200 e 0.14T1 / 2 or = e 0.14T1 / 2  t1/2 = - ln 0.5   4.81 years
2 0.14 0.14

Problems for discussion by the faculty in the class room

Session 12: Modeling and Solution of exponential growth and decay phenomena problems
1. A bacteria culture initially contains 100 cells and grows at a rate proportional to its size. After
an hour the population has increased to 420. a) Find an expression for the number of bacteria after
t hours. b) Find the number of bacteria after 3 hours. c) When will the population reach 10,000?
D)Find the rate of growth after 3 hours.

29
2. A radioactive material has an initial mass 100mg. After two years it is left to 75mg. Find the
amount of the material at any time. What is the period of its half-life?

3.If in a reactor Uranium losses 10% of its weight within one day, what is the half life? How
long would it take for 99% of the original amount to disappear?
Home Work
1. If the growth rate of the number of bacteria at any time t is proportional to the
number of bacteria present at t and doubles in 1 week, how many bacteria can be
expected after 2 weeks? After 4 weeks?
2. The half-life of a certain radioactive material is 85 days. An initial amount of the
material has a mass of 801 kg. Write an exponential function that models the
decay of this material remains after 10 days.
3. A radioactive isotope has an initial mass 400mg, which two years later is 100mg.
Then determine the value k of the decay problem.
Session 13
Newton’s law of cooling:
Newton’s law of cooling states that the rate of change of temperature of a cooling body is
proportional to the difference between the temperature of the body and that of the surrounding
medium.
Suppose that a body whose temperature is initially To is allowed to cool in air which
is maintained at a constant temperature of TS . It is required to find the temperature of the body as
functions of time .

Let the temperature of the body be at time . Then by Newton’s law of cooling

dT
 k (T  TS ) (1)
dt
where is the constant of proportionality.

Examples
1. A thermometer, reading 50C, is brought in a room whose temperature is 220C. One minute later
the thermometer reading is 120C.How long does it take until the reading is practically 22 0C, say
21.90C?
Sol:

30
 According to Newton's law f cooling, we have
   0  Ae  kt  t  0,   50 C and  0  220 C  A  17
10
when t  1 ,   120 C  k  log
17
when   220 C , t  ?  t  9.7 min .

2.If the temperature of a cake is 3000F when it leaves the oven and is 2000F ten minutes
later, when will it be practically equal to the room temperature of 60 0F say, when will it be
610F ?
Sol:
According to Newton's law f cooling, we have
   0  Ae kt  t  0,   3000 F and  0  600 F  A  240
1 24
when t  10 ,   2000 F  k  log
10 14
when   61 F , t  ?
0
 t  1hr 43min .
Problems for discussion by the faculty in the class room
Session -13: Modeling and solution of the Newton’s law of cooling.
1. A body kept in air with temperature 250C cools from 1400 C to 800 C in 20 min. Obtain
the time period, when the body cools down to 350 C.
2. A metal ball is heated to a temperature of 100oC and at time𝑡 = 0, it is placed in water
which is maintained at 40°C. If the temperature of the body is reduced to 60°C in 4 minutes,
find the time at which the temperature of the ball is 50°C.
3. If the temperature of a cake is 3000F when it leaves the oven and is 2000F ten minutes
later, when will it be practically equal to the room temperature of 60 0F say, when will it
be 610F ?
4. A cup of coffee (temperature = 190°F) is placed in a room whose temperature is 70°F.
After five minutes, the temperature of the coffee has dropped to 160°F. How many more
minutes must elapse before the temperature of the coffee is 130°F?

Home Work

1. A thermometer, reading 50C, is brought in a room whose temperature is 22 0C. One minute
later the thermometer reading is 120C.How long does it take until the reading is practically
220C, say 21.90C?
2. The body of a murder victim was discovered at 8.00pm. The doctor took the temperature
of the body at 8.30pm. Which was 94.6°F. He again took the temperature after one hour
when it was showed 93.4°F, and noticed that the temperature of the room was 70°F.
Estimate the time of death. (Normal temperature of human body 98.6°F).

31
 3. A hard-boiled egg at 70°C is put in a pan under running 20°C water to cool. After 5
minutes, the egg’s temperature is found to be 50°C. How much longer will it take the egg
to reach 30°C
Linear differential equations of higher order with constant coefficient

Linear Differential Equations are those in which the dependent variable and its derivatives occur
only in the first degree and are not multiplied together. The general linear differential equation
of the nth order is of the form

+ 𝑎 + 𝑎 + … + 𝑎 𝑦 = 𝑓(𝑥) (1)

Where 𝑎 , 𝑎 , . . . ,𝑎 and f are functions of 𝑥 or constants.

If 𝑎 , 𝑎 , . . . ,𝑎 are constants, then the above equation is called linear differential equation with
constant coefficients. Such equation are most important in the study of elecro mechanical
vibrations and other engineering problems

If 𝑦 , 𝑦 , . . . , 𝑦 are the linearly independent solutions of +𝑎 + … +𝑎 𝑦 =0,

where 𝑎 , 𝑎 , . . . ,𝑎 are constants, then 𝑐 𝑦 + 𝑐 𝑦 + ---- + 𝑐 𝑦 is also solution.

Solution of homogeneous equations

Consider the nth order homogeneous differential equation with constant coefficients

+𝑘 +𝑘 + . . . + 𝑘 y = 0, where 𝑘 , 𝑘 , . . . , 𝑘 are real constants.

Denoting , ,..., by D, D2, ----- D n, then we can write the above equation in symbolic
form as 𝐷 𝑦 +𝑘 𝐷 𝑦+𝑘 𝐷 𝑦 + . . . + 𝑘 𝑦 = 0  f(D)y=0,
where 𝑓 (𝐷) = 𝐷 +𝑘 𝐷 +𝑘 𝐷 + . . . + 𝑘 , the equation
𝑓 (𝑟) = 𝑟 + 𝑘 𝑟 + 𝑘 𝑟 + . . . + 𝑘 = 0 is called the Auxiliary Equation (A.E).
Let 𝑟 , 𝑟 , . . . , 𝑟 be its roots. Depending up on the nature of these roots, we can find the general
solution (complementary function) of 𝑓 (𝐷)𝑦 = 0, from the following table.
Roots of A.E C.F
1.𝑟 , 𝑟 , 𝑟 , . . . are real and different roots 𝑐 𝑒 +𝑐 𝑒 +𝑐 𝑒 +...

2.𝑟 , r1 , 𝑟 , . . . (two real and equal roots ) (𝑐 𝑥 + 𝑐 ) 𝑒 +𝑐 𝑒 +...

(𝑐 𝑥2+ 𝑐 𝑥 + 𝑐 ) 𝑒 +𝑐 𝑒 +...
3 r1 , r1 , r1 , r4 . . . (three real and equal roots)

4.𝛼 + 𝑖 , 𝛼 − 𝑖 , 𝑟 , . . . 𝑒 ∝ (𝑐 𝑐𝑜𝑠 𝑥 +𝑐 𝑠𝑖𝑛 𝑥) + 𝑐 𝑒 +...

32
 (a pair imaginary roots)
5.𝛼 ± 𝑖 , 𝛼 ± 𝑖 , 𝑟 , . . . 𝑒 ∝ [(𝑐 𝑥 + 𝑐 ) 𝑐𝑜𝑠 𝑥 + (𝑐 𝑥 + 𝑐 ) 𝑠𝑖𝑛 𝑥 ]
(two pairs of equal imaginary roots) +𝑐 𝑒 +...

Example 1 Solve 𝑦′′′ – 7𝑦′′ + 14 𝑦′ – 8y = 0.

Solution Given D.E in symbolic form is (𝐷3 – 7𝐷 2 + 14𝐷 −8)𝑦 = 0 (1)
Auxiliary equation of (1) is r3 – 7𝑟2 + 14𝑟 − 8 = 0.
Observe that 𝑚 = 1 is a root of this equation.
By synthetic division, we have 1 1 −7 14 −8
0 1 −6 8
1 −6 8 0
 (𝑟 − 1) (𝑟2 – 6𝑟 + 8) = 0
 (𝑟 − 1) (𝑟 − 2) (𝑟 − 4) = 0
The roots are 𝑟 = 1, 2, 4. Since roots are real and distinct, the general solution of the given
equation is 𝑦(x) = 𝑐 𝑒 + 𝑐 𝑒 + 𝑐 𝑒 .
Example 2 Solve the initial value problem 4𝑦 + 4𝑦 + 𝑦 = 0, 𝑦 (0) = 3, 𝑦 (0) = −7/2.
Solution Writing given D.E in symbolic form we get (4𝐷2 + 4𝐷 + 1) 𝑦 = 0
Auxiliary equation is 4𝑟2 + 4𝑟 + 1 = 0  (2r + 1)2 = 0
Therefore roots are 𝑟 = −1/ , −1/ .

Since roots are real and equal, the general solution is

/
𝑦 (𝑥) = (𝑐 𝑥 + 𝑐 )𝑒 , (1)
where 𝑐 , 𝑐 are arbitrary constants. Now to find the constants c1, c2, we use the initial conditions
𝑦(0) = 3 and 𝑦′ (0) = −7/2.
Differtiating (1) , we have
/ / /
𝑦′(𝑥) = − ½ (𝑐 𝑥 + 𝑐 )𝑒 +𝑐 𝑒 = [𝑐 (2 - x) 𝑐 ]𝑒 (2)

Putting 𝑥 = 0 and 𝑦 (0) = 3 in (1), we have 𝑦 (0) = 𝑐  𝑐 = 3

33
 1
Putting 𝑥 = 0 and 𝑦 (0) = − in (2), we have − =  [c1 ( 2  0)  c 2 ]
2
 2c1  c2  7  c1  2 − = 𝑐 – = 𝑐 – = 𝑐 = – = −2

Therefore, solution of the given initial value problem is 𝑦 (𝑥) = (−2𝑥 + 3) 𝑒

Example 3 Solve 𝑦 – 2𝑦 + 5𝑦 = 0, given 𝑦 (0) = −3, 𝑦′ (0) = 1.
Solution Given equation in symbolic form is (𝐷2 – 2𝐷 + 5) 𝑦 = 0.
±√
Auxiliary equation is 𝑟2 – 2𝑟 + 5 = 0  r= = 1 ± 2𝑖.

Since roots are imaginary, the general solution is

𝑦 (𝑥) = 𝑒 ( 𝑐 𝑐𝑜𝑠 2𝑥 + 𝑐 𝑠𝑖𝑛 2𝑥) (1)
𝑦′ (𝑥) = 𝑒 (−2𝑐 𝑠𝑖𝑛 2𝑥 + 2 𝑐 𝑐𝑜𝑠 2𝑥 ) + 𝑐 𝑐𝑜𝑠 2𝑥 + 𝑐 𝑠𝑖𝑛 2𝑥 (2)
Putting 𝑥 = 0 , 𝑦 = −3 in (1), we get 𝑐 = −3
Putting 𝑥 = 0 , 𝑦 = 1 in (2), we get 1 = 2𝑐 − 3  c2 = 2
Substituting 𝑐 = −3, 𝑐 = 2 in (1), we have
𝑦 (𝑥) = 𝑒 (−3 𝑐𝑜𝑠 2𝑥 + 2 𝑠𝑖𝑛 2𝑥). which is the required solution.
Session: 14
Solution of higher order homogeneous ODE with constant coefficients

Problems for discussion by the faculty in the class room

1. Determine the charge on the capacitor in an LRC series circuit at t = 0.01 s when inductance 1
H, resistance 5, capacitance (1/6) F, Electromotive force is 0 V, q (0) =3 C, and i (0) =0 A.
2.Determine charge q and current i in the LRC circuit with inductance 1H,resistance 8 ohms,
capacitance (1/16)F, E(t)=0, and the initial conditions are q(0)=0, i(0)=1.
3. Determine the charge on the capacitor in an LC series circuit at t when L = 2 h, C= 0.005 F, E
(t) = 0 V, q(0) =5 C, and i(0) =0 A.
Home Work

d2y dy
1. Solve the 3  2 y  0 , given that y(0) = 0, y (0)  1
2 dx
dx
2. Determine the charge on the capacitor in an LRC series circuit at t when inductance 1 H,
resistance 2 and capacitance 1/37 F, E(t)= 0 V, q(0) =2 C, and i(0) =0 A.
3. Determine the charge on the capacitor in an LRC series circuit at t when inductance 1 H,
resistance 4, capacitance 0.25 F, E(t)= 0 V, q(0) =5 C, and i(0) =0 A.

Session: 15

34
Solution of non – homogeneous equations

The nth order non–homogeneous with constant coefficient is

+𝑘 + ⋯ + 𝑘 𝑦 = 𝐹 (𝑥),

whose symbolic form is (𝐷n + 𝑘 𝐷n-1 + . . . + 𝑘 ) y = F(x) .

The general solution of the non - homogeneous equation is 𝑦 (𝑥) = 𝑦 (𝑥) + 𝑦 (𝑥), where

𝑦 (𝑥) : complementary function (or) general solution of the corresponding homogeneous

equation.
𝑦 (𝑥) : particular integral (or) particular solution of the non-homogeneous equation.

Rules for finding particular integral

F (x) is the P.I of the differential equation 𝑓(𝐷) 𝑦 = 𝐹(𝑥).

( )

Rules for finding the P.I: Consider the D.E 𝑓 (𝐷)𝑦 = 𝐹 (𝑥)
Case 1: 𝐹(𝑥) = 𝑒

i. If 𝑓 (𝑎)  0, then P.I = 𝑒 = 𝑒

( ) ( )
ii. If 𝑓(𝑎) = 0 and 𝑓 (𝐷) = (𝐷 − 𝑎)m 𝑔 (𝑑) , 𝑔 (𝑎)  0, then

P.I = 𝑒 =( ) ( )
𝑒 = 𝑒 (or) 𝑒 , (𝑓 (m)(𝑎)  0)
( ) ! ( ) ( )

Case 2: 𝐹(𝑥) = 𝑠𝑖𝑛 (𝑎𝑥 + 𝑏) 𝑜𝑟 𝑐𝑜𝑠 (𝑎𝑥 + 𝑏)

i) If 𝑓 (−𝑎2) 0, then
P.I = [ 𝑠𝑖𝑛 (𝑎𝑥 + 𝑏) 𝑜𝑟 𝑐𝑜𝑠 (𝑎𝑥 + 𝑏) ] = [ 𝑠𝑖𝑛 (𝑎𝑥 + 𝑏) 𝑜𝑟 𝑐𝑜𝑠 (𝑎𝑥 + 𝑏)]
( ) ( )

ii) If 𝑓(−𝑎2) = 0, then

P.I = [ 𝑠𝑖𝑛 (𝑎𝑥 + 𝑏) 𝑜𝑟 𝑐𝑜𝑠 (𝑎𝑥 + 𝑏) ]

( )

= ( )(
[ (𝑎𝑥 + 𝑏) 𝑜𝑟 𝑐𝑜𝑠 (𝑎𝑥 + 𝑏) , provided 𝑓 (m) (−𝑎2)  0, 𝑚 = 1, 2, 3, …
)

Case 3: 𝐹(𝑥) = 𝑥 , 𝑚 being a positive integer. Then P.I. = 𝑥 = [𝑓(𝐷 )] 𝑥 .

( )

Expand [𝑓(𝐷)] in ascending powers of 𝐷 and operate on 𝑥 term by term.

35
 i.e., P.I = [ 𝑎 + 𝑎 𝐷 + 𝑎 𝐷2 + . . . + 𝑎 𝐷m] 𝑥 .
Since the (m + 1) th and higher derivatives of 𝑥 are zero, no need to consider the terms beyond
𝐷 m.
Case 4: 𝐹 (x) = 𝑒 𝑣 (𝑥) , 𝑣 (𝑥) is function of 𝑥.

P.I = 𝑒 𝑣 (𝑥) = 𝑒 𝑣 (𝑥)

( ) ( )

Case 5: 𝐹 (x) = xV , where V is any function of 𝑥.

( )
P.I = (𝑥𝑉) = 𝑥 (𝑉) − (𝑉)
( ) ( ) [ ( )]

Example 1 Find the Particular Integral of (𝐷2 + 5𝐷 + 6) 𝑦 = 𝑒

Solution The particular integral is P.I = 𝑒 = ( )

𝑒 =

𝑦 (𝑥) = .

Example 2 Find the Particular Integral of (𝐷2 + 𝐷 − 2) 𝑦 = 2 𝑠𝑖𝑛 ℎ𝑥.

Solution The particular integral is P.I = 2 𝑠𝑖𝑛 ℎ𝑥= (𝑒 – 𝑒 )

= 𝑒 - 𝑒 = ( )(
𝑒 - 𝑒 = 𝑒 − 𝑒
) ( )

= + 𝑒 . 𝑦 (𝑥) = + 𝑒 .

Example 3 Find the Particular Integral of (𝐷3 + 1) 𝑦 = 𝑐𝑜𝑠 (2𝑥 – 1).

Solution The Particular Integral is given by P.I. = 𝑐𝑜𝑠 (2𝑥 − 1)

Substituting 𝐷2 = −22 = −4 in the above equation, we get

P.I = 𝑐𝑜𝑠 (2𝑥 – 1) X 𝑐𝑜𝑠 (2𝑥 − 1) = 𝑐𝑜𝑠 (2𝑥 − 1)

( ) ( ( ))
= [ 𝑐𝑜𝑠 (2𝑥 − 1) – 8 𝑠𝑖𝑛 (2𝑥 – 1)].
( )

𝑦 (𝑥) = [ 𝑐𝑜𝑠 (2𝑥 − 1) – 8 𝑠𝑖𝑛 (2𝑥 – 1)]

Student Outcome: After completion of this topic, student able to model the second order ODE
and solving that ODE.
Problems for discussion by the faculty in the class room
1. Determine the charge on the capacitor in an LRC series circuit at t = 0.01 s when inductance
1 H, resistance 13, capacitance 0.025 F, E(t)= e2t V, q(0) =5 C, and i(0) =0 A.

36
 2. Determine charge q and current i in the LRC circuit with inductance 0.5H, resistance 10
ohms, capacitance 0.02F, , and the initial conditions are q(0)=0, i(0)=1, E(t)=e 5t +4

3. Determine the charge on the capacitor in an LRC series circuit at t when inductance 1 H,
resistance 10, capacitance 0.04 F, E (t) = cos3t V, q(0) =5 C, and i(0) =0 A.

Home Work
1. Determine charge q and current i in the LRC circuit with inductance 2H, resistance 36
ohms, capacitance (1/144)F, E(t)=2+e2t, and the initial conditions are q(0)=0, i(0)=1.
2. Determine charge q and current i in the LRC circuit with inductance 0.5H, resistance 6
ohms, capacitance (1/16)F, E(t)=sinht , and the initial conditions are q(0)=0, i(0)=1.
3. Determine the steady-state charge and the steady-state current in an LRC series circuit
when inductance 1 H, resistance 2 ohms capacitance 0.25 F, and E (t) =50 cos (t+5) V.

Session16

Vibrations of a spring

Airplanes, bridges, ships, machines, cars etc., are vibrating mechanical systems. The simplest
mechanical vibrating system is mass – spring mechanical system. It consists of a coil spring of
natural length L, suspended vertically from a fixed point support. A constant mass ‘𝑚’ attached
to the lower end as the spring, stretches the spring to a length L + e and comes to rest, which is
known as the static equilibrium position. Here 𝑒 > 0 is the static elongation due to the hanging
the mass on the spring. Now the mass is set in motion either by pushing or pulling the mass from
equilibrium position and (or) by imparting a non zero velocity to the mass. Since the motion takes
place in the vertical direction. We consider the downward direction as positive.

Unstressed
spring e

System in
static 𝑥(𝑡)
equilibrium
System in motion

37
Let 𝑥 (𝑡) be the displacement of mass from the static equilibrium position. In order to determine
𝑥 (𝑡), we use Newton’s second law and Hooke’s law of motion. The mass ‘𝑚’ is subjected to the
following forces:
1. Gravitational force mg acting downwards.
2. Spring restoring force – 𝑘 (𝑥 (𝑡) + 𝑒) due to displacement of the spring from the rest.
3. Frictional force of the medium, opposing the motion and of magnitude – 𝑐 𝑥 ′(𝑡).
4. External force f (t).

By Newton’s second law

Mass  acceleration = total forces acting on the body

 𝑚 = 𝑚𝑔 – 𝑘 (𝑥 + 𝑒) – 𝑐 + 𝑓 (𝑡)

when 𝑘 > 0 is known spring constant or stiffness of the spring, 𝑐 ≥ 0 is known as damping
constant, g is gravitational constant.
Since the force on the mass exerted by the spring must be equal and opposite to the gravitational
force on the mass, we have 𝑘𝑒 = 𝑚𝑔.

Therefore 𝑚 = – 𝑘𝑥– 𝑐 + 𝑓 (𝑡)  𝑚 𝑥′′ + 𝑐 𝑥′ + 𝑘𝑥 = 𝑓 (𝑡)

Which is the differential equation describing the motion of the mass spring.
Session 16 Applications of second order ODE
Problems for discussion by the faculty in the class room
1. The motion of a mass on a certain vertical spring is described by
d 2x dx
 10  21x  e 2t  5, with x(0 )  0 and x (0)  2 , where x is the distance of the mass from
2 dt
dt
the equilibrium position, downward being taken as positive direction. Determine the displacement
of the motion.
2. The motion of a mass on a certain vertical spring is described by
d 2x dx
 12  32 x  e (5t  3) , with x(0 )  1 and x (0)  2 , where x is the distance of the mass
2 dt
dt
from the equilibrium position, downward being taken as positive direction. Determine the
displacement of the motion.

Home Work

1. The motion of a mass on a certain vertical spring is described by

10 y   100 y   90 y  0, y  0   0.16, y   0   0 , where y is the distance of the mass from the
equilibrium position, downward being taken as positive direction. Determine the
displacement of the motion
38
 2. The motion of a mass on a certain vertical spring is described by
x  16 x  64 x  0, x  0   0.33, x  0   0 , where x is the distance of the mass from the
equilibrium position, downward being taken as positive direction. Determine the
displacement of the motion.

COURSE OUTCOME 2
Session 17

Laplace Transforms
LEARNING OBJECTIVES:
After reading this chapter, the student will be able to understand:
 Introduction to Laplace transforms.
 Properties of Laplace transforms.
 Inverse Laplace Transforms.
 Convolution Theorem.
 Application of Laplace Transforms to solutions of Ordinary Differential Equations,
Simultaneous differential equations.

4.1 Introduction:
Many mathematical problems are solved using transformations. The idea is to transform the
problem into another problem that is easier to solve. Once a solution is obtained, the inverse
transform is used to obtain the solution to the original problem. The Laplace transform is an
important tool that makes solution of linear constant coefficient differential equations much easier.
The Laplace transform transforms the differential equations into algebraic equations which are
easier to manipulate and solve. Once the solution is obtained in the Laplace transform domain is
obtained, the inverse transform is used to obtain the solution to the differential equation. Laplace
transform is an essential tool for the study of linear time-invariant systems.

Laplace transform is yet another operational tool for solving constant coefficients linear
differential equations. The process of solution consists of three main steps:
 The given “hard” problem is transformed into a “simple” equation.
 This simple equation is solved by purely algebraic manipulations.
 The solution of the simple equation is transformed back to obtain the solution of the given
problem.
In this way the Laplace transformation reduces the problem of solving a differential equation to an
algebraic problem. The third step is made easier by tables, whose role is similar to that of integral
tables in integration.
The above procedure can be summarized by Figure.

39
The theory of Laplace Transforms is an essential part of the mathematical background required by
engineers, physicists and mathematicians. It gives an easy and effective means for solving certain
types of differential and integral equations. It is the foundation of the modern form of operational
calculus, which was originated in an attempt to justify certain operational methods used by an
electrical engineer Oliver Heaviside, in the later part of the 19 th century for solving equations in
electromagnetic theory.
Laplace Transforms help in solving the differential equations with boundary values without
finding the general solution and the values of the arbitrary constants.
Definition:

Let 𝑓 (𝑡) be a function defined for all positive values of 𝑡, then 𝐿[𝑓(𝑡)] = ∫ 𝑒 f (t) dt, provided
the integral value exists, is called the Laplace transform of f (t).
It is denoted as 𝐿 [𝑓 (𝑡)] = 𝐹 (𝑠) = ∫ 𝑒 𝑓 (𝑡)𝑑𝑡 , where‘s’ being a parameter.
Some important Formulae:

(1) L [1 ] =
!
(2) L [𝑡 ]= when n=0, 1, 2, 3…
(3) L [𝑒 ] = ( )
(𝑠 > 𝑎)
(4) L [sin 𝑎𝑡] = (𝑠 > 0)
(5) L [𝑐𝑜𝑠𝑡 𝑎𝑡 ] = (𝑠 > 0)
(6) L [sinh 𝑎𝑡] = (𝑠 > 𝑎 )
(7) L [𝑐𝑜𝑠ℎ 𝑎𝑡] = (𝑠 > 𝑎 )

Proofs of above Formulae:

(1) We know that L [ f(t)] = F (s) = ∫ 𝑒 f (t) dt

∴ 𝐿 [1 ] = ∫ 𝑒 𝑑𝑡 = = = [0 − 1] =
1
∴ 𝐿 [1] = .
𝑠
(2) We know that L [ f(t)] = F (s) = ∫ 𝑒 f (t) dt

∴ 𝐿[𝑡 ] = 𝑒 𝑡 𝑑𝑡 = 𝑒 𝑡 𝑑𝑡

40
 put 𝑠𝑡 = 𝑥 ⟹ 𝑠 𝑑𝑡 = 𝑑𝑥 (or) t= ∴ 𝑑𝑡 =

When t=0 ⟹ 𝑥 = 0 and 𝑡 = ∞ ⟹ 𝑥 = ∞

!
=∫ 𝑒 = ∫ 𝑒 𝑥 𝑑𝑥 = (Γ(𝑛 + 1) =

𝑛!
∴ 𝐿 [𝑡 ] = 𝑤ℎ𝑒𝑟𝑒 𝑛 𝑎𝑛𝑑 𝑠 𝑎𝑟𝑒 𝑝𝑜𝑠𝑡𝑖𝑣𝑒.
𝑠
Note: If n is a positive integer, then (n+1) = ∫ 𝑒 𝑥 𝑑𝑥 and (n+1) =𝑛!.

(3) We know that L [ f(t)] = F (s) = ∫ 𝑒 f (t) dt

( )
𝑒 ( )
∴ 𝐿 [𝑒 ] = 𝑒 𝑒 𝑑𝑡 = 𝑒 𝑑𝑡 =
−(𝑠 − 𝑎)

= { ( ) = ( )
{0 − 1} = .
( ) ( )

∴ 𝐿 [𝑒 ] = ( )
.

(4)We know that L [f (t)] = F (s) = ∫ 𝑒 f (t) dt

∵ 𝐿 [sin 𝑎𝑡 ] = ∫ 𝑒 sin 𝑎𝑡 𝑑𝑡 = ∫ 𝑒 𝑑𝑡

∴ sin 𝑎𝑡 = = ∫ 𝑒 𝑒 𝑑𝑡 − ∫ 𝑒 𝑒 𝑑𝑡

1
= 𝐿 𝑒 −𝐿 𝑒
2𝑖

= ( )
−( )
= =
( )
𝑎
∴ 𝐿 [sin 𝑎𝑡 ] =
(𝑠 + 𝑎 )
Similarly L [𝑐𝑜𝑠𝑎𝑡 ] = ∵ cos 𝑎𝑡 =
( )

(5) We know that L [f (t)] = F (s) = ∫ 𝑒 f (t) dt

∴ 𝐿 [sinh 𝑎𝑡 ] = 𝑒 sinh 𝑎𝑡 𝑑𝑡
𝑒 −𝑒
= 𝑒 𝑑𝑡
2
∵ sin 𝑎𝑡 = and cos at =

= ∫ 𝑒 𝑒 𝑑𝑡 − ∫ 𝑒 𝑒 𝑑𝑡

41
= {𝐿 [𝑒 ] − 𝐿 [𝑒 ]}

= − =
( ) ( ) ( )
𝑎
∴ 𝐿 [sinh 𝑎𝑡 ] =
(𝑠 − 𝑎 )

Similarly L [cosh 𝑎𝑡 ] =
( )

4.2 Properties of Laplace transform:

4.2.1 Linear Property:

If L [f (t)] = F (s) = ∫ 𝑒 f (t) dt, then L [𝑎𝑓 (𝑡) + 𝑏𝑓 (𝑡)] = 𝑎. 𝐿 𝑓 ( ) + 𝑏. 𝐿 [𝑓 (𝑡)]

Proof: -We know that L [f (t)] = F (s) = ∫ 𝑒 f (t) dt.

∴ 𝐿[𝑎𝑓 (𝑡) + 𝑏𝑓 (𝑡)] = 𝑒 [𝑎𝑓 (𝑡) + 𝑏𝑓 (𝑡)]𝑑𝑡

=𝑎 ∫ 𝑒 𝑓 (𝑡)𝑑𝑡 + 𝑏. ∫ 𝑒 𝑓 (𝑡)𝑑𝑡

= 𝑎. 𝐿 [𝑓 (𝑡)] + 𝑏. 𝐿 [𝑓 (𝑡)]
∴ L[𝑎𝑓 (𝑡) + 𝑏𝑓 (𝑡)] = 𝑎. 𝐿 [𝑓 (𝑡)] + 𝑏. 𝐿 [𝑓 (𝑡)].
4.2.2 First Shifting property:

If L [𝑓(𝑡)] = 𝐹 (𝑠) = ∫ 𝑒 f (t) dt, then L [𝑒 𝑓(𝑡)] = 𝐹(𝑠 − 𝑎).

Proof:-We know that L [𝑓(𝑡)] = 𝐹 (𝑠) = ∫ 𝑒 f (t) dt

∴ 𝐿 [𝑒 𝑓 (𝑡)] = ∫ 𝑒 𝑒 𝑓 (𝑡)𝑑𝑡
( )
= ∫ 𝑒 f (t) dt.
𝑝𝑢𝑡 (𝑠 − 𝑎) = 𝑟
=∫ 𝑒 𝑓 (𝑡)𝑑𝑡 = 𝐹 (𝑟) = 𝐹 (𝑠 − 𝑎).
∴ 𝐿 [𝑒 𝑓 (𝑡)] = 𝐹 (𝑠 − 𝑎)
With the help of this property, we can have the following important results:

!
1. L[𝑒 𝑡 ] = ( )

2. 𝐿 [𝑒 sin 𝑏𝑡] =
( )
( )
3. 𝐿 [𝑒 cos 𝑏𝑡] = ( )

4. 𝐿 [𝑒 sinh 𝑏𝑡] = ( )
( )
5. 𝐿 [𝑒 cosh 𝑏𝑡] = ( )

42
Change of scale property:

If L [𝑓 (𝑡)] = 𝐹 (𝑠), 𝑡ℎ𝑒𝑛 𝐿 [𝑓(𝑎𝑡)]= 𝐹

Proof: -We know that L [𝑓 (𝑡)] = 𝐹 (𝑠) = ∫ 𝑒 f (t) dt
∴ 𝐿 [𝑓(𝑎𝑡)] = ∫ 𝑒 𝑓(𝑎𝑡)𝑑𝑡

𝑝𝑢𝑡 𝑎𝑡 = 𝑢 ∴ 𝑑𝑡 =

𝐿[𝑓(𝑎𝑡)] == ∫ 𝑒 𝑓(u) = ∫ 𝑒 f (u) 𝑑𝑢 = 𝐹 .

Second shifting property:

𝑓(𝑡 − 𝑎), 𝑡 > 𝑎
If L [𝑓 (𝑡)] = 𝐹 (𝑠) 𝑎𝑛𝑑 𝑔 (𝑡) = 𝑡ℎ𝑒𝑛 𝐿 [𝑔(𝑡)] = 𝑒 𝐹 (𝑠).
0 𝑡<𝑎
Proof: -We know that L [𝑓 (𝑡)] = 𝐹 (𝑠) = ∫ 𝑒 f (t) dt

∴ 𝐿 [𝑔(𝑡)] = 𝑒 𝑔 (𝑡)𝑑𝑡

=∫ 𝑒 𝑔 (𝑡)𝑑𝑡 + ∫ 𝑒 𝑔(𝑡)𝑑𝑡
= 0 +∫ 𝑒 𝑓 (𝑡 − 𝑎)𝑑𝑡
𝑝𝑢𝑡 (𝑡 − 𝑎) = 𝑢 𝑑𝑡 = 𝑑𝑢
𝑤ℎ𝑒𝑛 𝑡 = 𝑎 ⟹ 𝑢 = 0 𝑎𝑛𝑑 𝑡∞ ⟹ 𝑢∞ .
( )
∴ 𝐿 [𝑔 (𝑡)] = 𝑒 𝑓 (𝑢)𝑑𝑢

( )
=∫ 𝑒 . 𝑒 𝑓 (𝑢)𝑑𝑢 = 𝑒 =𝑒 . 𝐹 (𝑠).

Example 1. Find the Laplace Transform of co𝑠 𝑡 .

Solution. 𝐿 [𝑐𝑜𝑠 𝑡] = 𝐿 = 𝐿 [1 + 𝑐𝑜𝑠2𝑡]

= {𝐿 [1] + 𝐿 [𝑐𝑜𝑠2𝑡] [∵ 𝑐𝑜𝑠2𝑡 = 2𝑐𝑜𝑠 𝑡 − 1]

= + . [∵ 𝑐𝑜𝑠 𝑡 = ]

Example 2. Find the Laplace Transform of 𝑡 .

! ( )
Solution. We know that 𝐿 [𝑡 ] = 𝑜𝑟

43
 √⊼ ⊼
Put 𝑛 = ∴ 𝐿 𝑡 = = = = .
√

Example 3.Find the Laplace transform of 𝑡. 𝑠𝑖𝑛 𝑎𝑡 .

Solution. 𝐿 [𝑡. 𝑠𝑖𝑛𝑎𝑡] = 𝐿 𝑡. = 𝐿 𝑡. 𝑒 −𝑒

= 𝐿𝑒 𝑡 −𝐿 𝑒 𝑡

= − = .
( ) ( ) ( )

Example 4. Find the Laplace Transform of 𝑒 (2𝑐𝑜𝑠5𝑡 − 3𝑠𝑖𝑛5𝑡).

Solution. We know that 𝐿 [2𝑐𝑜𝑠5𝑡 − 3𝑠𝑖𝑛5𝑡]=𝐿[2𝑐𝑜𝑠5𝑡] − 𝐿[3𝑠𝑖𝑛5𝑡]

=2 −3 =
( )
Therefore, 𝐿 [ 𝑒 (2𝑐𝑜𝑠5𝑡 − 3𝑠𝑖𝑛5𝑡)]= = . [By using shifting property]
( )

Session 17: Laplace Transforms

Problems for discussion by the faculty in the class room

1. Determine L  (2t  5) 2 

2. Determine L  e 2 t sinh t 

3. Determine L  t sin 2t 

4. Apply first shifting property and Determine L  t 2 e 2t 

Home Work
1. Determine the following :
(a) 𝐿 [𝑒 + 3𝑒 ].
(b) 𝐿 [3𝑒 + 5𝑐𝑜𝑠𝑡].
(c) 𝐿 [sinh 6𝑡 + 3 𝑒 + 𝑐𝑜𝑠5𝑡].
(d) 𝐿 [𝑠𝑖𝑛 2𝑡].
(e) 𝐿 [𝑠𝑖𝑛3𝑡. 𝑠𝑖𝑛2𝑡].
(f) 𝐿 [𝑐𝑜𝑠4𝑡. 𝑠𝑖𝑛2𝑡].
(g) 𝐿 [(𝑡 + 1) ].
(h) 𝐿 [5𝑒 + 𝑐𝑜𝑠ℎ3𝑡 + 𝑠𝑖𝑛5𝑡].
(i) 𝐿 [𝑒 − 𝑒 ].
( )
(j) 𝐿 𝑒

44
Session 18
Inverse Laplace Transforms:
Now, we obtain f (t) when F (s) is given, then we say that inverse Laplace Transform of F(s) is f
(t). If L [𝑓(𝑡)] = 𝐹 (𝑠), 𝑡ℎ𝑒𝑛 𝐿 [𝐹(𝑠)] =
𝑓(𝑡), 𝑤ℎ𝑒𝑟𝑒 𝐿 𝑖𝑠 𝑐𝑎𝑙𝑙𝑒𝑑 𝑡ℎ𝑒 𝑖𝑛𝑣𝑒𝑟𝑠𝑒 𝐿𝑎𝑝𝑙𝑎𝑐𝑒 𝑇𝑟𝑎𝑛𝑠𝑓𝑜𝑟𝑚 𝑜𝑝𝑒𝑟𝑎𝑡𝑜𝑟.
From the application point of view, the Inverse Laplace Transform is very useful.
Important Formulae:

(1) 𝐿 =1

(2) 𝐿 =( )

(3) 𝐿 =𝑒
(4) 𝐿 = cos 𝑎𝑡.
(5) 𝐿 = sin 𝑎𝑡
(6) 𝐿 = cosh 𝑎𝑡
(7) 𝐿 = sinh 𝑎𝑡
(8) 𝐿 [𝐹 (𝑠 − 𝑎)] = 𝑒 𝑓(𝑡)
(9) 𝐿 ( )
= 𝑒 𝑠𝑖𝑛𝑏𝑡
( )
(10) 𝐿 ( )
= 𝑒 𝑐𝑜𝑠𝑏𝑡

(11) 𝐿 ( )
= 𝑒 𝑠𝑖𝑛𝑏𝑡
( )
(12) 𝐿 ( )
= 𝑒 𝑐𝑜𝑠𝑏𝑡

(13) 𝐿 ( )
= 𝑡 sin 𝑎𝑡

(14) 𝐿 ( )
= 𝑡 cos 𝑎𝑡

(15) 𝐿 ( )
= [sin 𝑎𝑡 + 𝑎𝑡 cos 𝑎𝑡]

45
Note - Partial Fractions:

Denominator Expression Form of Partial Fractions

containing…
a. Linear factor

b. Repeated linear factors

c. Quadratic term
(which cannot be
factored)

Example 1. Find the Inverse Laplace Transform of the following:

(a) ( )
(b) ( )
(c) ( )
(d) ( )

( )
(e) ( )
(f) ( )
(g) ( )
(h) ( )

Solution.

(a) 𝐿 ( )
= 𝑒 .𝐿 =𝑒 1=𝑒 .

(b) 𝐿 ( )
=𝐿 = 𝐿 = sinh 3𝑡.

(c) 𝐿 ( )
=𝐿 ( )
= cosh 4𝑡.

(d) 𝐿 ( )
=𝐿 ( )
= 𝐿 = sin 5𝑡.

(e) 𝐿 ( )
=𝑒 𝐿 ( )
= 𝑒 . sin 2𝑡.

( ) ( )
(f) 𝐿 ( )
=𝐿 ( )
=𝑒 𝐿 = 𝑒 cos 2𝑡.

(g) 𝐿 ( )
=𝑒 𝐿 =𝑒 . 𝐿 = 𝑒 .sinh2t.

46
(h) 𝐿 =𝐿 = 𝐿 = 𝑒 .

Some Important Theorems on Inverse Laplace Transforms:

(a) First Shifting theorem:𝐼𝑓 𝐿 [𝐹 (𝑠)]=f (t) then 𝐿 [𝐹 (𝑠 − 𝑎)] = 𝑒 𝑓(𝑡)
(b) Second Shifting theorem:𝐼𝑓 𝐿 [𝐹 (𝑠)]=f (t);
𝑓(𝑡 − 𝑎) , 𝑡 > 𝑎
𝐿 [𝑒 𝐹(𝑠)] = 𝐺 (𝑡) 𝑤ℎ𝑒𝑟𝑒 𝐺 (𝑡) =
0 𝑡<𝑎

(c) Change of Scale Property :𝐼𝑓 𝐿 [𝐹 (𝑠)]=f (t) then 𝐿 [𝐹 (𝑎𝑠)] = 𝑓 ,𝑎 > 0 ,

(d) Inverse Laplace Transform of Derivatives :𝐼𝑓 𝐿 [𝐹 (𝑠)]=f (t) then

( )
𝐿 = (−1)𝑡 𝑓(𝑡).

(e) Inverse Laplace Transform of Integrals : 𝐼𝑓 𝐿 [𝐹 (𝑠)]=f (t) then

( )
𝐿 ∫ 𝐹(𝑠)𝑑𝑠 =

(f) Multiplication by power of s:𝐼𝑓 L [F (s)]=f (t) and f (0) =0 then 𝐿 [𝑠. 𝐹(𝑠)]=𝑓 (𝑡)

(g) Division by powers of s:

( )
If f (t) is Piecewise continuous and of exponential order ‘a’ and lim exists then for s>
→
( )
𝑎, 𝐿 = ∫ 𝑓(𝑥)𝑑𝑥 .

Example 1. Find the Inverse Laplace Transform of the following:

( ) ( )
(a) 𝐿 =𝐿 ( )
=𝐿 ( )
−𝐿
( )

=𝑒 . cos 3𝑡 − 2. 𝑒 . sin 3𝑡 =𝑒 cos 3𝑡 − sin 3𝑡 .

(𝑏)𝐿 =𝐿 ( )
= 𝐿 = 𝐿

= 𝑒 𝐿 = 𝑒 .𝑡

(c) 𝐿 ( )
= 𝐿𝑒𝑡 𝐹 (𝑠) = ( )
⟹ [𝐹(𝑠)] = ( )

∴ 𝐹 (𝑠) = ∫ ( )
𝑑𝑠 =
( )

𝑤𝑒 𝑘𝑛𝑜𝑤 𝐿 [𝐹 (𝑠)] = (−𝑡). 𝐿 [𝐹(𝑠)] = (−t). 𝐿 ( )

= .𝐿 ( )
= .𝑠𝑖𝑛𝑎𝑡.

47
 ( )
(d). 𝐿 ( )
=∫ 𝐿 𝑑𝑡 = ∫ . 𝑠𝑖𝑛𝑎𝑡 𝑑𝑡 = {−𝑐𝑜𝑠𝑎𝑡} =

(e) 𝐿 𝑙𝑜𝑔 =

𝑙𝑒𝑡 𝐿 𝑙𝑜𝑔 =f(t)

1+𝑠 𝑑 1+𝑠
∴ 𝐿 [𝑓(𝑡)] = 𝑙𝑜𝑔 , 𝑡ℎ𝑒𝑛 𝐿 [𝑡. 𝑓(𝑡)] = (−1) 𝑙𝑜𝑔
𝑠 𝑑𝑠 𝑠
𝑑 2 1
= (−1) [𝑙𝑜𝑔 (1 + 𝑠) − 𝑙𝑜𝑔𝑠 ] = −
𝑑𝑠 𝑠 (𝑠 + 1)
2 1
∵ [𝑡. 𝑓(𝑡)] = 𝐿 − = 2(1) − 𝑒 (1) = (2 − 𝑒 )
𝑠 (𝑠 + 1)

∴ [𝑡. 𝑓(𝑡)] = 𝐿 𝑙𝑜𝑔 = (2 − 𝑒 )

(f) 𝐿 =𝐿 ( )
−( )
=𝐿 ( )
−𝐿 ( )
=𝑒 −𝑒 .

Convolution Theorem:
Definition:- Let f(t) and g(t) be two functions of class A then the convolution of f(t) and g(t)
is denoted by f(t) * g(t) and is defined by f(t) * g(t) =∫ 𝑓(𝑥)𝑔(𝑡 − 𝑥)𝑑𝑥

Also f (t) * g (t) is called the resultant or falling of f(t) and g (t) .
Convolution Theorem: If f (t) and g(t) are two functions of class A (t ≥ 0) and 𝐿 [𝐹 (𝑠)]=f
(t), 𝐿 [𝐺 (𝑠)]=g (t) then

𝐿 [𝐹 (𝑠). 𝐺 (𝑠)]=∫ 𝑓(𝑥)𝑔(𝑡 − 𝑥)𝑑𝑥 = 𝑓(𝑡) ∗ 𝑔(𝑡).

Session 18
Problems for discussion by the faculty in the class room
 1 1 s 
L1    2
1. Find  s  3 s s  4 

 2s  3 
L1  2
2. Find  s  4 s  13 

 1 
L1  
3. Find   s  1 s  3 

s2
4. By Convolution theorem find L1 ( )
(s 2  a 2 ) (s 2  b2 )

48
Home Work
1. Find the inverse Laplace transform of ( )(
by Convolution theorem.
)

2. Find the inverse Laplace transform of ( )

by Convolution theorem.
3. Using Convolution theorem, find the Laplace transform of ( )
.

4. Find 𝐿 ( )
using Convolution theorem.

Session 19
Application of Laplace Transforms to solutions of Ordinary Differential Equations:
Laplace transformation is useful in solving ordinary as well as partial differential equations.
The advantage of this method over the conventional one, as already pointed out, is that it
reduces the problem of solving a differential equation to an algebraic problem. More over this
method yields the particular integral directly without finding complementary function and
particular integral and then evaluating the arbitrary constants.
Let us consider a linear differential equation with constant coefficients.

+𝑎 +𝑎 +-------+𝑎 + 𝑎 𝑦 = 𝑓(𝑡) -----------(1)

And y(0) =𝛼 , 𝑦 (0) = 𝛼 , -------- 𝑦 (0) = 𝛼 be the given initial or boundary conditions
where 𝛼 , 𝛼 , − − − − − − −𝛼 , are constants.
Taking Laplace Transforms of both sides of the differential equation with the help of the
formula
𝐿 [𝑦 (𝑡)]=𝑠 𝐿[𝑦] − 𝑠 𝑦(0) − 𝑠 𝑦 (0)--------𝑦 (0)
On using the given initial conditions, the equation (1) reduces to an algebraic equation of the
form y=F(s), which is called subsidiary equation, from we can find the required solution by
taking inverse Laplace transform of F(s). ie.,
𝐿 [𝑦]=𝐿 [𝐹(𝑠)] or 𝑦 = 𝑓(𝑡) 𝑖𝑠 𝑡ℎ𝑒 𝑟𝑒𝑞𝑢𝑖𝑟𝑒𝑑 𝑠𝑜𝑙𝑢𝑡𝑖𝑜𝑛.
Example 1 Solve by Laplace Transform method. 𝑦 − 3𝑦 + 2𝑦 = 4,
𝑤ℎ𝑒𝑟𝑒 𝑦(0) = 2; 𝑦 (0) = 3.
Solution. Given that 𝑦 − 3𝑦 + 2𝑦 = 4 ---------------(1)
Taking Laplace transform of both sides 𝐿 [𝑦 ] − 3. 𝐿[𝑦 ] + 2. 𝐿 [𝑦] = 4. 𝐿 [1]

{𝑠 . 𝐿[𝑦] − 𝑠. 𝑦(0)} − 3{𝑠. 𝐿[𝑦] − 𝑦(0)} + 2. 𝐿 [𝑦] = 4

using 𝑦(0) = 2, 𝑦 (0) = 3

[𝑠 . 𝐿[𝑦] − 𝑠. 2 − 3] − 3 [𝑠. 𝐿[𝑦] − 2] + 2 𝐿[𝑦]

= (𝑠 − 3𝑠 + 2)𝐿 [𝑦] = − 2𝑠 − 3 ∴ 𝐿 [𝑦] = ( )

.

49
 = −( )
+( )
(resolving into partial fractions).

∴ 𝑦 = 2. 𝐿 − 3. 𝐿 + 3. 𝐿 = 2 − 3𝑒 + 3𝑒 .

Example 2.Solve (𝐷 − 𝐷 − 2)𝑦 = 20 𝑠𝑖𝑛2𝑡 𝑤ℎ𝑒𝑟𝑒 𝑦(0) = 1, 𝑦 (0) = 2

Solution. Given that (𝐷 − 𝐷 − 2)𝑦 = 20 -----------(1)
Applying Laplace transforms on both sides
𝐿 [𝑦 (𝑡)] − 𝐿 [𝑦 (𝑡)] − 2𝐿[𝑦 (𝑡)] − 20. 𝐿 [𝑠𝑖𝑛2𝑡]

{𝑠 . 𝐿[𝑦] − 𝑠. 𝑦(0) − 𝑦 (0)} − {𝑠. 𝐿[𝑦] − 𝑦[0]} − 2𝐿[𝑦] = 20.

Using y (0) =1, 𝑦 (0)=2, we have

(𝑠 − 𝑠 − 2) 𝐿[𝑦] − 𝑠. 1 − 2 + 1 = .

∴ 𝐿[𝑦]= + 𝑠. 1 . =( )( )
.

𝐿 [𝑦] = − ( )
+ +
( )

= 𝐿 + 𝐿 +𝐿 − 3𝐿

= 𝑒 + 𝑒 + 𝑐𝑜𝑠2𝑡 − 3𝑠𝑖𝑛2𝑡.

Example 3 An 8-lb weight is hung on the end of a vertically suspended spring thereby
stretching the spring 6 inches. The weight is raised 3 inches above its equilibrium position and
released from rest at time t = 0. Find the displacement x of the weight from its equilibrium
position at time t. Use g= 32 ft/s2.

1
Solution: By Hooke’s law 8  k    k  8 lb/ft.
2

d 2x k
By Newton's second law  x  0, where x is measured positively downward.
dt 2 m
1
 From x  t   64 x  t   0, x  0    , x  0   0
4
Applying Laplace transform on both sides, we obtain

s 2 L  x   sx  0   x  0   64 L  x   0
1 s
 L  x  
4 s 2  64
1
Thus x  t    cos 8t
4

50
Session 19
Problems for discussion by the faculty in the class room

1. Solve y  y  0 given that y  0   1 , y   0   0 using Laplace transform

2. Solve y"+4y' +3y = et with initial conditions y(0) = 0, y'(0) = 2

3. The motion of a mass on a certain vertical spring is described by

10 y   100 y   90 y  0, y  0   0.16, y   0   0 , where y is the distance of the mass from the
equilibrium position, downward being taken as positive direction. Determine the
displacement of the motion by using Laplace transform.

Home Work

1. A particle moves along a line so that its displacement x from a fixed point O at any time
t is
given by 𝑥 + 4𝑥′ + 5𝑥 = 80𝑠𝑖𝑛5𝑡, if the particle starts from rest initially, find its
displacement at any time, t>0.
2. A resistance R in series with inductance 𝐿 is connected with e.m.f 𝐸(𝑡). The
current 𝑖(𝑡) in time 𝑡 is given by 𝐿 + Ri = E(t). If the switch is connected at
𝑡 = 0 and disconnected at 𝑡 = 𝑎, find the current 𝑖(𝑡).
3. Solve the following differential equation using Laplace transform
𝑦 + 2𝑦 + 2𝑦 = 5 sin 𝑡 , 𝑦(0) = 𝑦 (0) = 0.

Session 20

FOURIER SERIES
1. Introduction
In many physical and engineering problems, particularly those connected with
vibrations and conduction of heat, it is more useful to be able to express a real valued function
in series of sines and co-sines. Most of the single-valued functions which occur in applied
mathematics can be expressed in the form.
a1+a2 Cos x +a2 Cos 2x+……. + b1 Sin x + b2 Sin 2x +……
within a desired range of values of the variable. Such a series is known as
Fourier Series.
The use of Fourier series in special problems dated from the time of Daniel Bernoulli
(1700-1782) who used these to solve certain problems connected with vibrating strings. The
systematic study of the subject was, however, first undertaken by the French mathematician
Jacques Fourier (1768-1833) in his memorable monograph “Theorie Analytique de la chaleur”.
2. Periodic Functions
If the function f(x) satisfies for all values of x the relation
f (x+T) = f(x),

51
 where T is a real number, the function is said to be periodic if T is the smallest positive
number for which such a relation is satisfied then T is called period of the function. The graph
of such a function is obtained by periodic repetition of its graph in any interval of length T.
when the above relation is satisfied, we find that
f(x) = f(x+T) = f (x+2T)
= ……… = f (x+nT) =…..
and also f(x) = f (x-T)
= f (x-2T) = ……
= f (x-nT) = ……
i.e., f (x) = f (x + nT),

where n is positive integer. Thus f(x) repeats itself again and again after periods T.
The simplest periodic functions are Sin x and Cos x, having the period 2  . Their
reciprocals Cos ec x and Sec x are also periodic with period 2  ; Tan x and Cot x are periodic
with period  .
The functions Sin nx and Cos nx are periodic with period 2  /n. The functions.
2x 2x
Sin and Cos
T T
are periodic with period T.
The sum of a number of functions with commensurable periods will also be a periodic
function. For example
Sin x + 1
2 Sin 2x + 1
3 Sin 3x + ……. ….(i)

is a periodic function. The periods of the various terms are 2  ,  , 23  respectively. So the
period of function (i) is the least common multiple of these periods, viz., 2  .
In the subsequent articles we shall use the following results.
a  2 a  2
 1 
i) a Sin nx dx =  n Cos nx a =0 (n  0)

a  2 a  2
 1 
ii) a Cos nx dx   n Sin nx a =0 (n  0)

a  2 a  2
iii)  Sin nx Cos nx dx 
a
1
2  Sin (m  n) x  Sin (m+n)xdx=0
a

a  2 a  2
1
iv)  Sin mx Sin nx dx 
a
2  [Cos (m  n) x  Cos (m  n) x] dx
a

52
 =0 (m  n)
a  2 a  2
1
v)  Cos mx Cos nx dx 
a
2  [Cos (m  n) x  Cos (m  n) x] dx
a

=0 (m  n)
a  2 a 2 x
1
 Sin nx dx   (1  Cos 2nx) dx   (n  0)
2
vi)
a
2 a

a  2 a  2
1
a Cos nx dx  2  (1  Cos 2nx) dx   (n  0)
2
vii)
a

Euler’s Formulae
It can be shown that under certain conditions a function f(x) can be expressed as n
infinite series of sines and cosines of x and its integral multiples. Thus if f(x) be a function of
x known in interval c<x<  +c we can write

a 0*
f (x) = + a1 Cos x+ a2 Cos 2x + ….an Cos nx + ….
2
= b1 Sin x + b2 Sin 2x + ….+ bn Sin nx+ …. ….(1)
a0
* to write instead of a0 is a conventional device to be able to get more symmetric.
2

a0
i.e., f(x) =
2
 (a
n 1
n Cos nx  bn Sin x)

c  2
1
where a0 =
 
c
f ( x ) dx …(2)

c  2
1
an =
 c
f ( x) Cos nx dx ….(3)

c  2
1
bn =
 c
f ( x) Sin nx dx …(4)

These values of a0, an, bn are known as Euler’s formulae and the series (1) as Fourier
series for f(x).
Proof. Let f(x) be represented in the interval (c, c+2x) by the Fourier series.

a0
f(x) =
2
  a
n 1
n Cos nx bn Sin nx  ...(5)

To find the coefficients a , an, bn we assume that the series (5) can be integrated term
by term from x=c to x =c+2  .
To find a0 integrating both sides of (5) from x = c to x= c+2  , we have

53
 c  2 c  2 c  2
1   
 f ( x) dx  a 0  dx +    a n Cos nx  dx
c
2 c c  n 1 
c  2
  
+    bn Sin nx  dx
c  n 1 
(by integrals I and II of Art. 2)
1
= a 0 (c  2  c) + 0 + 0
2
= a0 
c  2
1
Hence a0 =
  f ( x) dx
c

To find a n , multiplying each side of (5) by Cos nx and integrating from x= c to x=c
+2  , we get
c  2 c  2
1

c
f ( x) Cos nx dx =
2
a  Cos nx dx
c

c  2
  
+    a n Cos nx  Cos nx dx
c  n 1 
c  2
  
+    bn sin nx  Cos nx dx
c  n 1 

= +  an  0

(by integrals I, III, IV, V and VI)

c  2
1
Hence an =
  f ( x) Cos nx dx
c

To find bn, multiplying each side of (5) by Sin nx and integrating from x=c to x=c+2
 , we get
c  2 c  2
1

c
f ( x) Sin nx dx=
2
a  Sin nx dx
c

c  2
  
+    a n Cos nx  Sin nx dx
c  n 1 
c  2
  
+    bn Sin nx  Sin nx dx
c  n 1 

= 0+   bn

54
 (by integrals II, V, VI, VII, & VIII)
c  2
1
Hence bn=
  f ( x) Sin nx dx
c

Corollary, 1 making c=0, the interval becomes 0<x<2  and the formulae (2), (3) and (4)
reduce to
2
1
a0 
  f ( x) dx
0
(6)

2
1
an
  f ( x) Cos nx dx
0
(7)

2
1
bn 
  f ( x) Sin nx dx
0
(8)

Corollary 2.Putting c=-  , the interval becomes -  >x<  and the formulae (2), (3) and (4)
reduce to

1
a0 
  f ( x) dx

(9)


1
an
  f ( x) Cos nx dx

(10)


1
bn 
  f ( x) Sin nx dx

(11)

Conditions for a Fourier Expansion

The reader must not be misled by the belief that the Fourier expansion of f(x) in each
case shall be valid. The above discussion has merely shown that if f(x) has an expansion, then
the coefficients are given by Euler’s formulae. The problems concerning the possibility of
express in g a function by Fourier series and the convergence are many and cumbersome. Such
questions should be left to the curiosity of a pure-mathematician. However, almost all
engineering applications are covered by the following conditions known as Dirichlet’s
conditions.
Any function f(x) can be expanded as a Fourier Series

a0
  ( a n Cosnx  bn Sinnx)
2 n 1
Provided
(i) f(x) is periodic, single-valued and finite;
(ii) f(x) has a finite number of finite discontinuities in any one period;
(iii) f(x) has at the most finite number of maximum and minimum.
These conditions are known as Dirichlet conditions.

55
For example:
1
(i) Let f ( x) 
x2
This function is discontinuous at x=2 but it is infinite discontinuity. Hence it does not
satisfy Dirichlet conditions.
1
(ii) Let f ( x)  Sin
x2
This function has infinite number of maximum and minimum. Hence it does not
satisfy Dirichlet conditions.
(iii) Step function. It is defined as
H(t) = 0, t<0
= 1, t  1
This function satisfies Dirichlet conditions.
ILLUSTRATIVE EXAMPLES
Example 1. Find the Fourier series to represent the function, f(x) given by
f(x)= -x,    x  0,
f(x)= x, 0<x< 

a0
Taking f(x)=   a n Cosnx  bn Sinnx
2 n 1
We have
 0 
1 1 1
a0 
 

f ( x)dx 
  ( x)dx   0
x dx

 
=  
2 2
 0 
1 1 1
an 
 

f ( x)Cos nx dx 
  ( x)Cox nx dx   0
x Cos nx dx

0 
1  xSin x Cosnx  1  xSin nx Cos nx 
=    +  
 n n     n
2
n 2  0

1 Cos n Cos n 1
=   
n 2
n 2
n 2
 n2
2
=  (Cos n  1).
 n2

56
  0 
1 1 1
bn 
  f ( x) Sin nx dx    ( x) Sinnx dx    xSin nx dx
  0

0 
1   Cos nx Sinnx  1  xCos nx Sin nx 
=   2  +  
 n n     n n 2  0

Cos n Cos n
=  0
n n
Consequently , we get the series
 4 1 1 
F(x)=   Cos x  2 Cos 3 x  2 Cos 5 x  ....
2  3 5 
Observation. Putting x=0 in (ii), we find another interesting series involving
 4 1 1 1 
0  1  2  2  2  ....
2  3 5 7 

2 1 1 1
i..e, 1 2
 2  2  .....
8 3 5 7
Example 2. Find a Fourier series to represent x-x2 from x= -  to x= 
Taking f(x)=x-x2

a0
   ( a n Cos nx b n Sin nx) (i)
2 n 1
Then

1 2 2
a0   ( x  x ) dx  
2

 
3

1
an   ( x  x
2
)Cos nx dx
 


1 Sinnx  Cos nx   Sin x 
 ( x  x ) 2  (1  2 x)     ( 2) 3 
 n  n 2
  n   

4
 Cosnx
n2
4 4 4 4
 a1  2
, a 2   2 , a 3  3 , a 4  2 etc
1 2 3 4

1
bn   ( x  x
2
)Sin nx dx
 

57
 
1  Cos nx   Sinnx   Cos nx  
 ( x  x 2 )     (1  2 x) 2   ( 2) 
  n   n   n
3
  

2
 Cosn
n
 b1  2, b2  1, b3  2 / 3, b4  2 / 4 etc.

Hence

2  Cosx Cos 2 x Cos3x Cos 4 x 

(x  x2 )    4 2  2
 2
 2
 ......
3  1 2 3 4 

 Sinx Sin 2 x Sin3 x Sin 4 x 

+2      ......
 1 2 3 4 
Observation. Where x=0, , we find
2 1 1 1 1
 2
 2  2  2  .....
12 1 2 3 4
Example 3. Obtain the Fourier Series for ex in the interval -  <x< 
Taking f(x)= ex

a0
   ( a n cCos nx b n Sin nx) ………….(1)
2 n 1
We have
 
1 1
a0 
2 
 f ( x)dx  
2 
e x dx

1 
 (e  e  )
2
 
1 1
an  f ( x)Cos nx dx   e Cos nx dx
x

 
 


 (n
1
2
 1)
e Cosnx  ne
x 
Sin x 



 (n
1
2
 1)
e Cosnx  ne Cos nx
x x

 
1 1
bn   f ( x) Sin nx dx   e
x
Sin nx dx
 
 


 (n
1
2
 1)
e Sinnx  ne Cos nx 
x x 


58
 
n
 ( n 2  1)

e  Cosnx  ne  Cos nx 
 n Cos nx 
 ( e  e  )
 ( n  1)
2

Subsequently, we get the series

e   e   1  1 1 1  1 2 3 
ex     Cosx  Cos 2 x  Cos 3 x....   Sin x  Sin 2 x  Sin3 x.....
 2  2 5 10  2 5 10 
Example 4. Obtain Fourier series expansion for f(x) defined as follows

f ( x)  x  ,  x  0
2

f ( x)   x,0  x  
2
We find here
 
1 1 
0
 1 
a0 
 

f ( x)dx    x  dx   (  x)dx
   2  0 2

1  x 2   x2  
0

    x   x   
  2 2    2 2  0 
 

1  2 2 2 2
     0
 2 2 2 2 


1
an 
  f ( x)Cos nx dx



1 
0
 1  
   x  Cos nx dx     x Sinnx dx
   2  0 2 

1   Sinnx Cos nx  
0
  Sinnx Cos nx   
  x        x    
  2 n n 2    2  n n 2  0 

1  1 Cos nx Cos nx 1 
    2
  n 2 n2 n2 n 
2
 (1  Cosnx )
 n2

59
 4
where n is odd
  n 2
0 when is even


1
bn 
  f ( x)Sin nx dx



1 
0
 1  
   x  Sin nx dx     x Sinnx dx
   2  0 2 

1    Cos nx Sin nx  
0
  Cos nx Sin nx   
    x         x    
   2 n n 2     2  n n 2  0 

1     
    Cos nx  Cos nx    0
  2n 2n 2n 2n 
Subsequently, the series is
41 1 1 
f ( x)   Cos x  2 Cos 3 x  2 Cos 5 x  .................
 12
3 5 
Notes:1 In evaluating the coefficients in a Fourier Series we often have to use integration by
parts repeatedly. It is, therefore, convenient to obtain a standard for it. Thus

 ur dx  ur 1  u r 2 u r3  u r4....................................

2.Fourier series, being a sum of periodic functions, is itself a periodic function. So a function
f(x) will be represented faithfully by its Fourier series within the specified interval only. Beyond
this interval the sum will be repetition of the values for the specified interval.
3. In example 1 the function f(x) has been defined by two different formulae in two
different regions. Such types of functions though rarely used in ordinary analysis,
are fairly common in Fourier series and Engineering problems.
Session 20
1. Determine the Fourier series to represent the function f(x)=(π-x)from x= -π to x = π.
2. Construct the Periodic function for the graph in the interval (0, 2π),

(π, π)

(0, 0) (π,0) (2π, 0)

and then express it as a Fourier Series.
3. Identify the periodic function from the following wave form in the interval (-π, π).

π
π
60
-4π
Home Work

1. Determine bn in the Fourier series to represent the function e-x from x= 0 to x = 2π.
2. Find a0 and an in the Fourier expansion for the function f(x) = x-x2 ; -1<x <1.
3. Draw the graph for the function f(x) is defined as follows and express as Fourier
series.

Session 21
Even and Odd Functions

(i) A function f(x) is said to be odd, if f(x-)= -f(x).

e.g., Sin x Tan x, x2 all odd functions. Graphically , an odd function is symmetrical about the
origin.
(ii) A function f(x) is said to be even, if f(-x)= f(x).

e.g., Cos x, Sin x2, x2 all even functions. Graphically , an even function is symmetrical about
the y-axis.

We shall be using the following property of definite integrals in the next paragraph.
c

 f ( x)dx  0, when f ( x) is an odd

c
function , .


 2 f ( x ) dx, when f ( x) is an even function , .
c

Expansion of Odd or Even Functions

We know that a periodic function f(x) defined in ( ,  ) can be represented by the Fourier series

a0 nx nx
f ( x) 
2
 a
n 1
n cos
c
 bn sin
c
c
1
where a0  
c c
f ( x ) dx

nx
c
1
a n   f ( x) Cos dx
c c c
nx
c
1
bn   f ( x) Sin dx
c c c
Case I. When f(x) is an odd function.

61
 c
1
a0  
c c
f ( x ) dx  0

nx nx
Since Cos is an even function, therefore f(x) Cos is an odd function.
c c
nx
c
1
 a n   f ( x ) Cos dx =0
c c c
nx nx
Again, Since Sin is an odd function, therefore f(x) .Sin is an even function.
c c

c
1 nx
 bn =
c  f ( x)
c
Sin
c
dx

nx
c
2
=
c 0 f ( x) Sin
c
dx

Hence, if a periodic function f(x) is odd, its Fourier expansion contains only Sin e terms and

c
2 nx
bn =
c  f ( x) Sin
0 c
dx .

Case II. When f(x) is an even function.

c c
1 2
a0  
c c
f ( x) dx   f ( x) dx
c0
nx
Since f(x) Cos is also an even function
c
nx
c
1
 a n   f ( x) Cos dx .
c c c
nx
c
2
=
c  f ( x) Cos
0
c
dx

nx
Again, Since f(x) Sin is an odd function
c
c
1 nx
 bn = 
c c
f ( x ) Sin
c
dx =0

Thus if a periodic function f(x) is even, its Fourier expansion contains only cosine terms,
and

c
2
a0 =
c  f ( x) dx
0
c
2 nx
a0 =
c  f ( x) Cos
0 c
dx.

62
Example 1. Find a Fourier series to represent x2 in the interval (-l, l)

Since f(x) = x2, is an even function in (-l l)

a0 
nx
 f(x) =
2
  an Cos
n 1 l
..(i)

l
2 2l 2
 x dx 
2
Then a0 =
l 0
3

l
2 nx
x
2
an = Cos dx
l 0 l
l
2  2  Sin nx / l   Cos nx / l   Sin nx / l 
=
l  x  n / l   2 x   n 2  2 / l 2   2  n 3 3 / l 3 
       0

4l 2
= Cos n  .
n 2 2

 4l 2 4l 2  4l 2  4l 2
 a1 = , a2 = , a 3 = , a 4 = , etc.
2 22  2 32  2 42  2

Substituting these values in (i), we get

l 2 4l 2  Cos x / l Cos 2x / l Cos 3x / l 

x2 =     
3 2  12 22 32 

the required series.

Example 2: Express x as a Fourier series from   to 

The function is odd in the given interval.
Therefore, a0 = an = 0


2
also bn =
  f (x) Sin
0
nx dx


2
=
 x
0
Sin nx dx


2   Cos nx   Sin nx 
= x     
   n   n 2  0
2( 1) n
bn = -
n
Hence for    x  
x=2 (Sin x - 1
2 Sin 2 x  13 Sin 3 x  14 Sin 4 x  ...........) .

63
Observations

As the series has a period 2  , it represents the discontinuous function. It is important

to note that the given function y =x is continuous, but the function represented by the series (i)
has finite discontinuities at x=   , x=  3 , x=  5 etc.
Example 3. Find Fourier expansion for the function Sin ax in the interval -l<x<l
Here f(x) = Sin ax, -l<x<l
 f (-x) = - Sin ax = - f(x)
 The function is odd and hence
an = an =0,
nx
l 1
2 2 nax
and bn =  f ( x) Sin dx   Sin ax Sin dx
l 0 l l 0 l

 n  n
l
1    
= 
l 0Cos 
 l
 a  x  Cos 
  l
 a  x  dx
 
l
1 1  nx  1  n  
=  Sin   a  x  Sin   a  x
l  n  al  1  n  al  l  0
1 1
= Sin ( n  al )  Sin ( n  al )
n  al n  al
1 1 2
Thus, b1 = Sin al  Sin al  2 Sin al
  al   al   a 2l 2
1 1  4
b2 = Sin al  Sin al  Sin al
2  al 2  al 4  a 2 l 2
2

1 1 6
b3 = Sin al  Sin al  Sin al etc
3  al 3  al 9  a 2 l 2
2

Thus the series is

 1 x 2 2x 3 3x 
f(x) = 2  Sin al  2 Sin  2 2 Sin  2 2 Sin  ....
  a l l 2  a l 3  a l
2 2 2 2 2 2
l l 
Session 21
1. Compute the Fourier series for the function f(x) =|x| in– π x  π.

2. Compute the Fourier coefficients for the function f(x)=x3-x in -1<x<1 and f(x+2)=f(x).

Home Work

1. Obtain the Fourier series for x2 in the interval    x   .

1 1 1 1 2
Hence show that 1-    ....... 
2 2 32 4 2 52 12
2. Obtain Fourier series for the function f(x) given by
f(x) = 1+ 2 x /  ,    x  0
= 1- 2 x /  , 0<x< 

64
Session 22
Half Range Series

In some problems it is desired to expand a function in a Fourier series of a function f(x)

for the interval (0,c) which is half the period of the Fourier series. In addition, we may be forced
by the conditions of the problem to expand a given function in a series of sines alone or a series
of consines only. As it is immaterial whatever the function may be outside the range 0<x<c,
we extent the function to cover the range –c<x<c so that the new function may be odd or even.
The Fourier expansion of such a function of half the period, therefore, consists of Sine or cosine
terms only.

Sin e series: If it be required to expand f(x) as series of since in 0<x<c, then we extend the
function reflecting it in the origin, so that
F(-x)=-f(x)
Then extended function is odd in (-c,c) and its expansion will give the Fourier series :

nx
f ( x )   bn Sin
n 1 c
where
nx
c
2
bn   f ( x)Sin dx.
c0 c
Cosine series. If it be required to express f(x) as a cosine series in 0<x<c, we extend the
function reflecting it in the y-axis, so that
f(x) = +f(x)
Then the extended function is even in (-c, c) and its expansion will give fourier cosine series:
a 
nx
f(x) = 0   a n Cos
2 n 1 c
where
c
2
c 0
a0 = f ( x) dx

c
2 nx
and 
c 0
f ( x) Cos
an =
c
dx.

ILLUSTRATIVE EXAMPLES

Example1.
Find the Fourier series half-range Cosine series for the function f(x)=x, 0<x<  .


1
Here an 
  f ( x) dx
0

1 1

 xdx  
0
2

2

  x Cos nx dx
0

65
 
2  Sin nx  Cos nx  
 x   
 n  n 2   0
4
 0( n  even, 2 ( n  odd )
n
Hence the Fourier Cosine series is
 4 Cos 3 x Cos 5 x
  Cosx    .... 
2  32 52
1 1
Example 2. Expand f(x)=  x, if 0  x 
4 2
3 1
= x  , if  x  1
4 2

in the Fourier series of Sine terms

Let f(x) represent an odd function in (-1,1) so that


f ( x )   bn Sinnx
n 1

1
where bn  2 f ( x) Sin nx dx
0

 12 1

 1 3 
 2   (  x)Sinnx dx   ( x  ) Sinnx dx 
4 4
0 1

 2 
1
1
 1 Cos nx Sin nx  2   3  Cosnx Sinnx 
 2 (   x )  2 2   2   x    2 2 
 4 n n  0   4  n n  1
2

1 4 Sin ( n / 2)
bn  (1  Cosn )
2 n n 2 2
Thus
1 4   1 4   1 4 
f ( x )    2  Sinx    2 2  Sin3x    2 2  Sin5x  .........
    3 3    5 5  

Example 3. Express Sin x as a Cosine series in 0<x<  .


a.0
Let Sinx    a n Cosnx
2 n 1
 
2 2
where a0 
 
0
f ( x)dx 
  Sinxdx
0

2 4
=  Cosx 0 
 

66
 
2
and an 
  f ( x)Cos nx dx
0


2

  Sin x Cos nx dx
0


 Sin(n  1) x  Sin(n  1) xdx

1

 0

1  Cos(n  1) x Cos(n  1) x 
 
  n  1 n  1  0
1  Cos ( n  1) Cos ( n  1) 1 1 
    
 n 1 n 1 n  1 n  1

=0, when n is odd.

2 4  Cos 2 x Cos 4 x Cos 6 x 

Hence f ( x)       ........
   1. 2 3. 5 5.7 


Example 4: If f ( x)  ,0  x   / 3
3

= 0,  / 3  x  2 / 3

=   / 3,2 / 3  x  

2  1 1 
Then f ( x)  Cosx  Cos 5 x  Cos 7 x  ....
3 5 7 

1 1
and also f(x)=Sin 2x+ Sin 4 x  Sin10 x  .......... ..
2 10

(i) For Cosine Series


2
a0 
  f ( x)dx
0

 3 2


2  3
 
=   dx   0 dx    dx 
 03 3
  2

 3 3 

2   2  2 2 2 
=   0
  9 3 9 

67
 
2
an 
  f ( x)Cos nx dx
0

 2


2 3  3
 
   Cos nx dx   0.Cos nx dx    Cos nx dx 
 03 3
  2

 3 3 


  
2   Sinnx  3  Sinnx  
 .    
 3  n  0  n  2 
 3 

2 1 n 1 2 n 
  Sin  Sin
3 n 3 n 3 

2 n  n 
 Sin 1  2 Cos 3 
3n 3
Thus
2   
a1  Sin 1  2 Cos 3 
3 3

2 3 2
 . .(1  1) 
3 2 3

2 2  2 
a2  Sin 1  2Cos 3   0
3. 2 3

2  3 
a3  Sin 1  2Cos   0
3. 3  3 

2 4  4 
a4  Sin 1  2Cos 3   0
3. 4 3

2 5  5 
a5  Sin 1  2Cos 3 
3.5 3

2 3.2 2
 .  .etc.
3.5 2 5 3

2  1 1 
 f ( x)  Cosx  Cos 5 x  Cos 7 x.......
3 5 7 

68
(ii) For Sine Series


2
bn 
  f ( x) Sin nx dx
0
 2
3 

2  3
 
   Sin nx dx   0.Sinnx dx    Sin nx dx 
 03 3
  2

 3 3 


  
2   Cosnx  3  Cos nx  
 .      
 3  n 0  n  2 
 3 

2 1 n 1 1 1 2 n 
   Cos   Cos nx  Cos
3 n 3 n n n 3 

2 1 n 2 n 
   Cos  1  Cos n  Cos
3 n 3 3 

2   2 
Thus b1   Cos  1  Cos  Cos 0
3.1  3 3 
2  2 4 
b2    Cos  1  Cos 2  Cos 1
3 .2  3 3 
2  3 6 
b3    Cos  1  Cos 3  Cos 0
3.3  3 3 

2  4 8  1
b4   Cos  1  Cos 4  Cos   etc.
3.4  3 3  2

1 1
Hence f ( x )  Sin 2 x  Sin 4 x  Sin10 x  .........
2 10
Session 22

1. Find the Fourier series half-range Cosine series for the function f(x)=x, 0<x< 
2. Express Sin x as a Cosine series in 0<x<  .
3. Find the half-range Sine series for the function f ( x)  x (l  x) for 0  x  l

Home Work

1. Find the half-range Sine series for the following function

f ( x )  x (  x ) for 0  x  
2. Find the half-range Cosine series for the following function:
f ( x )  2 x  1 for 0  x  1

69
 3. Show that half range Sine series for f(x) is
8l 2 
1  x 
 3  (2n  1)
n 0
2
sin (2n  1)  where f ( x)  lx  x 2 in (0, l ),
 l 
Session 23
Partial differential equations

We will now consider differential equations that model change where there is more than one
independent variable. For example, the temperature in an object changes with time and with
the position within the object. The rates of change lead to partial derivatives, and the
equations relating them are called partial differential equations. The applications of the
subject are many, and the types of equations that arise have a great deal of variety. We will
limit our study to the equations that arise most frequently in applications. These model heat
flow and simple waves. solving them in some cases using the method of separation of
variables.

Partial differential equation is an equation relating to the partial derivatives of some unknown
functions.
Goals: 1. Formulate PDE, 2. Solve PDE
A partial differential equation (PDE) is an equation involving one or more partial derivatives
of an unknown function that depends on two or more variables, often time t and one or several
variables in space
z  2 z
Ex:  2
t x
The order of the highest derivative in PDE is called as the order of the PDE.
z  2 z
Ex:  2 second order PDE
t x
Standard notation: Let x and y represent the independent variables and z the dependent
variable so that 𝑧 = 𝑓 (𝑥, 𝑦) then
z z
 z x  p,  z y  q, and
x y
2 z 2 z 2z
 z xx  r ,  z xy  s ,  z yy  t.
x 2 xy y 2
These are the first , second order parial derivatives of z  f(x, y) respectively.
Just like as in the case of ordinary differential equations, second order PDE’s will be the most
important ones in applications.
A Partial differential equation can be obtained by elimination of arbitrary constants or by
elimination of arbitrary functions, involving two or more variables.
Here we are discussing about how to obtain partial differential equations by eliminating
arbitrary constants and functions from the solution.
Elimination of arbitrary constants:
Consider z to be a function of two independent variables x and y defined by
𝑓(𝑥, 𝑦, 𝑧, 𝑎, 𝑏) = 0 ---- (1) in which a and b are arbitrary constants.
To Eliminate two constants we need at least three equations .Therefore partially differentiating
above equation w.r.t x and y we get two or more equations. From these three equations we can
eliminate the two constants ‘a’ and ‘b’
Differentiating equation (1) partially with respect to x and y, we obtain

70
 𝜕𝑓 𝜕𝑓 𝜕𝑧 𝜕𝑓 𝜕𝑓
+ = +𝑝 =0 − −(2)
𝜕𝑥 𝜕𝑧 𝜕𝑥 𝜕𝑥 𝜕𝑧
𝜕𝑓 𝜕𝑓 𝜕𝑧 𝜕𝑓 𝜕𝑓
+ = +𝑞 =0 − −(3)
𝜕𝑦 𝜕𝑧 𝜕𝑦 𝜕𝑦 𝜕𝑧
By means of above three equations, two arbitrary constants can be eliminated.
So that we obtained a result i.e., a required Partial differential equation of order one in the form
𝐹(𝑥, 𝑦, 𝑧, 𝑝, 𝑞) = 0
Elimination of arbitrary functions:
Let 𝑢 = 𝑢(𝑥, 𝑦, 𝑧) 𝑎𝑛𝑑 𝑣 = 𝑣(𝑥, 𝑦, 𝑧)𝑏𝑒 𝑖𝑛𝑑𝑒𝑝𝑒𝑛𝑑𝑒𝑛𝑡 functions of the variables x, y ,z and
let ∅(𝑢, 𝑣) = 0 − (1) be an arbitrary relation between them . We shall obtain the Partial
Differential Equation by eliminating the functions 𝑢 and 𝑣. Regarding z as the dependent
variable and differentiating (1) partially with respect to x and y, then we get
  u u z    v v z 
     0
u  x z x  v  x z x 
  u u    v v 
  p    p  0
u  x z  v  x z 
and
  u u    v v 
  q     q   0
u  y z  v  y z 
e lim inating
 
, from(2)and (3) wehave
u v

u u v v
p p
x z x z
=0
u u v v
q q
y z y z
 u u   v v   u u   v v 
i.e.   p    q  -   q    p  =0
 x z   y z   y z   x z 
this is the required differential equation.
Writing 𝑃 = −
𝑄= −
𝑅= −
Then the above equation takes of the form 𝑃𝑝 + 𝑄𝑞 = 𝑅 a partial differential in p and q free
of the arbitrary function ∅(𝑢, 𝑣).Thus from this equation which involves arbitrary function ∅
we obtain a partial differential equation which is linear.
Example 1. Formulate the PDEs from the relation 2z=x2/a2+y2/b2 by eliminating the arbitrary
constants a, b.
Solution Given that 2z = x2/a2 + y2/b2 (1)
Differentiate equation (1) partially w.r.to x and y to obtain
1 1 z 1 1 z
 ; 
a 2
x x b 2 y y

71
Substituting these values in equation (1), we get the required PDE 2z=px+qy.
Example 2. Formulate the partial differential equation of all spheres whose centre lie on Z-
axis and given by equation x2 + y2 + (z-a)2 = b2, a and b being constants.
Solution. Equation of the sphere given by x2 + y2 + (z-a)2 = b2 (1)
Where a and b are two arbitrary constants.
Differentiating (1) partially w. r. t ‘x’ and ‘y’
x + (z-a) p=0 (2)
y + (z-a) q=0 (3)
Eliminating ‘a’ from equations (2) & (3), we get
yp - xq=0
This is the required partial differential equation.
Example 3. Formulate the partial differential equation from z=f (x2+y2).
Solution. Given z=f(x2+y2) (1)
Differentiating (1) partially w. r. t ‘x’ and ‘y’
z
p= = f ( x 2  y 2 )2 x (2)
x
z
q= = f ( x 2  y 2 )2 y (3)
y
p x
Dividing (2) by (3), we get   py  qx  0 which is required partial differential
q y
equation.
Example 4. Formulate a partial differential equation by eliminating the arbitrary function f
from the relation f ( x2 + y2, x2 - z2)=0.
Solution. The given equation f (x2 + y2, x2 - z2) = 0-----(1) can be expressed in the following
manner x2 + y2=g(x2 - z2) (2)
Differentiate equation (2) w.r.to x and y partially, we get
2x = g'(x2 - z2)(2x - 2zp) and 2y = g'(x2 - z2)(-2zq)
Dividing these two relations, we obtain the following relation
xz q + xy - zyp=0, which is a first order linear partial differential equation.
Session 23
Session Outcome: Formulate the partial differential equation from the solution and related
applications.

Problems to be discussed by the faculty in the class:

1. Formulate the partial differential equation of the following by eliminating the arbitrary
functions from z= 𝑓(𝑥 + 𝑐𝑡) + 𝑔(𝑥 − 𝑐𝑡)
1 
2. Formulate the partial differential equation of z  y 2  2 f   log y  by eliminating the
x 
arbitrary functions
3. Formulate the partial differential equation of all spheres of radius 5 and having their
centre on xy - plane
4. Formulate the partial differential equation of f ( x  y  z , x 2  y 2  z 2 )  0 by eliminating
the arbitrary Functions.
Homework Problems:
1. Formulate the partial differential equation of all spheres whose Centre lie on Z-axis and
given by Equation, 𝑥 + 𝑦 + (𝑧 − 𝑎) = 𝑏 a and b being constants.

72
2. Formulate the partial differential equation of the following by eliminating the arbitrary
Functions from z = 𝑓(𝑥 + 𝑦 + 𝑧 )
3. Formulate the Partial Differential Equation from the relation 2𝑧 = + by eliminating
the arbitrary constants a, b.
4. Obtain the differential equation of all planes cutting off equal intercepts with x and y axis
5. Formulate the partial differential equation of family of cones having vertex at the origin.

Session 24
Session outcome: Apply Lagrange Method to solve linear differential equations
Method of Grouping:
A linear PDE of the first order is of the form Pp + Qq = R (1)
Where P, Q and R are functions of x, y, z this equation is known as Lagrange’s linear PDE.
To obtain the solution of equation (1), we use the following steps:
dx dy dz
(1) Write the subsidiary equations in the form   .
P Q R
(2) Solve these simultaneous equations we obtain u(x, y) =c1 and v(x, y) =c2.
(3) Write the solution as f (u, v) =0 or u=f (v), or v=f (u).
Session 24
Method of Multipliers:
If the Lagrange Subsidiary equations are not exact, then choose one set of multipliers l 1, m1, n1
such that l1P + m1Q + n1R = 0 implies l1dx + m1dy + n1dz = 0 and then integrate, we get one
integral u (x, y) = c1.
Similarly choose another set of multipliers l2, m2, n2, then we adopt the above procedure, we
obtain another integral v (x, y) = c2.
Hence the solution is of the form f (u, v) = 0.
Example 1: Obtain the solution for pz – qz = z2 + (x+y)2.
dx dy dz
Solution The Lagrange’s Subsidiary equations are   2

z  z z  ( x  y) 2 
Consider the first two fractions and integrate we get one integral x + y = c 1.
Consider the first and last fraction and using the fact that x+y=c1 and then integrate, we obtain
the second integral log [z2 + (x + y)2] - 2x = c2.
Therefore, the general solution is f (x + y, log [z2 + (x + y)2] - 2x) = 0.
Example 2 Obtain the solution for y 2 p  xyq  x( z  2 y ).
Solution Given y 2 p  xyq  x( z  2 y )
dx dy dz
The Subsidary equations are 2   (1)
y  xy x( z  2 y )
dx dy
Consider first two members of the equations 
y x
On integration x 2  y 2  c1 (2)
dy dz
From last two members of equation (1), 
 y ( z  2 y)
On integration y  yz  c2
2
(3)
From (2) and (3) the required solution is f ( x 2  y 2 , y 2  yz ) =0.
Example 3 Obtain the solution for (mz - ny) p + (nx - lz) q = ly - mx.
dx dy dz
Solution The subsidiary equations are  
mz  ny nx  nz ly  mx

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 xdx  ydy  zdz
Using multipliers x, y, z we have, each fraction=  x2+y2+z2=c1.
0
Again consider another set of multipliers l, m, n we get ldx + mdy + ndz = 0 integrating, we
obtain second integral lx + my + nz = c2.
Hence the required solution of the give partial differential equation is
F (x2+y2+z2, lx + my + nz) = 0, where f is an arbitrary function.

Problems to be discussed by the faculty in the class:

1. Obtain the solution of partial differential equation 𝑝 tan 𝑥 + 𝑞 tan 𝑦 = tan 𝑧. Using
Lagrange’s method
2. Obtain the solution of partial differential equation 𝑥(𝑦 − 𝑧)𝑝 + 𝑦(𝑧 − 𝑥)𝑞 = 𝑧(𝑥 − 𝑦).
Using
Lagrange’s method
3. Obtain the solution of 𝑝√𝑥 + 𝑞 𝑦 = √𝑧 using Lagrange’s method
4. Obtain the solution of partial differential equation x 2 ( y  z ) p  y 2 ( z  x)q  z 2 ( x  y ) using
Lagrange’s method
Homework Problems:
1. Obtain the solution of partial differential equation p – q = log (x + y) using Lagrange’s
method
2. Obtain the solution of partial differential equation (y + z) p + (x + z) q = x + y using
Lagrange’s method
z z x 2  y 2
3. Obtain the solution of partial differential equation ( y  x)  ( y  x)  using
x y z
Lagrange’s method
4. Obtain the solution of partial differential equation
z z
x( y 2  z 2 )  y ( z 2  x 2 )  z ( x 2  y 2 )
x y

Session 25

Using Lagrange’s method

Many problems in engineering give rise to the following well-known partial differential
equations:
u  2u
(i) One dimensional heat flow equation:  c2 2 .
t x
2 y  2
y
(ii) Wave equation: 2  c 2 2 .
t x
 2u  2u
(iii) Two dimensional Laplace’s equation:  0.
x 2 y 2
There are many applications of partial differential equations frequently occur in the theory of
Elasticity and Hydraulics.
One of the methods to solve these equations is Method of separation of variables.
Session 25
Method of Separation of Variables:
In this method, we assume that dependent variable is the product of two functions each of
which involves one of the independent variables. So two ordinary differential equations are
formed.

74
 u u
Example 1 Using the method of separation of variables solve  2  u where
x t
u ( x,0)  6e 3 x .
u u
Solution Given differential equation is  2 u (1)
x t
Let 𝑢(𝑥, 𝑡) = 𝑋 𝑇 (2)
where 𝑋 is function of 𝑥 and 𝑇 is a function of 𝑡 only.
u u
 X 'T and  XT '
x t
X ' X T '
Substituting in equation (1) X 'T  2 XT ' X T and separating the variables 
2X T
[Where each is equal to some constant say K for all x and t]
X ' X T '
i.e.  K
2X T
X ' X
consider K
2X
on solving, we get X  c1e (1 2 K ) x .
T'
Now consider K.
T
On solving, we get T  c2 e kt
Hence, u ( x, t )  XT  c1e (1 2 K ) x c2 e kt  u ( x, t )  c1c2 e (1 2 K ) x e kt
Given 6e 3 x  u ( x,0)  c1c2 e (1 2 K ) x e kt
Comparing we get, c1c 2  6 and K  2
Therefore , u  6e  (3 x  2t ) is the required solution.
Problems to be discussed by the faculty in the class:
1. Apply method of Separation of variables to get the solution of the partial differential
equation = 4 = 0 u(0, y) = 8𝑒
2. Apply method of Separation of variables to get the solution of the partial differential
equation 3 − 2 = 0 u(x, 0) = 4𝑒
3. Apply method of Separation of variables to get the solution of the partial differential
equation 𝑢 = 𝑒 cos 𝑥 with 𝑢(𝑥, 0) = 0 𝑎𝑛𝑑 𝑢(0, 𝑡) = 0
Homework Problems:
1. Apply method of Separation of variables to get the solution of the partial differential
equation −2 + =0
2. Apply method of Separation of variables to get the solution of the partial differential
u u
equation x 2  y2 0
dx dy
3. Apply method of Separation of variables to get the solution of the partial differential
equation 2𝑥 − 3𝑦 = 0
4. Apply method of Separation of variables to get the solution of the partial differential
u u
equation  2  u where u ( x,0)  6e 3 x
x t

75
Session 26
Laplace Equation in Two Dimensions:
 2 u  2u
The differential equation of the form   0 is known as Laplace equation in two
x 2 y 2
dimensions.
Solution of Laplace Equation in two dimensions
 2 u  2u
Here the equation is  0
x 2 y 2
(1)
Let u  XY
where X is a function of x only and Y is a function of y only, be a solution of (1)
 2u  2u
Then  X ' ' Y and  XY ' '
x 2 y 2
Substituting in (1),we have
X '' Y ''
X ' 'Y  XY ' '  0 or 
X Y
(2)
Since x and Y are independent variables, this equation hold only when both sides reduce to a
constant say k. Equation 2 leads to ODE.
X ' 'kX  0 and Y ' ' kY  0
(3)
Solving equation (3), we get
i) When k is a positive and equal to p 2 (say)
X  c1e px  c 2 e  px , Y  c3 cos py  c 4 sin py
ii) When k is negative and equal to  p 2 (say)
X  c1 cos px  c 2 sin px , Y  c3 e py  c 4 e  py
iii) When k is equal to zero
X  c1 x  c2 , Y  c3 y  c4
Thus the various possible solutions of the Laplace equation are
u  (c1e px  c2 e  px )(c3 cos py  c 4 sin py )
u  (c1 cos px  c 2 sin px )(c3 e py  c 4 e  py )
u  (c1 x  c2 )(c3 y  c4 )
Of these three solutions we have to choose that solution which is consistent with the
physical nature of the problem and the given boundary conditions.
 2u  2 u
Example 1 Solve   0 which satisfies the conditions u(0,y)=u(l,y)=u(x,0)=0
x 2 y 2
nx
and u(x,a)= sin .
l
 2 u  2u
Solution The given equation is  0 (1)
x 2 y 2
The three possible solutions of (1) are
u  (c1e px  c2 e  px )(c3 cos py  c 4 sin py ) (2)

76
 u  (c1 cos px  c 2 sin px )(c3 e py  c 4 e  py ) (3)
u  (c1 x  c2 )(c3 y  c4 ) (4)
Keeping in view the given boundary conditions the only possible solution is (3)
Therefore, u ( x, y )  (c1 cos px  c2 sin px)(c3 e py  c 4 e  py ) (5)
Since u(0,y)=0  0  c1 (c3 e py  c 4 e  py )  c1  0
Equation (5) reduces to u ( x, y )  (c 2 sin px )(c3 e py
 c 4 e  py ) (6)
n
u(l,y)=0 0  (c2 sin pl )(c3 e py  c4 e  py )  sin pl  0  p 
l
nx
also u(x,0)=0  0  c2 sin (c3  c4 )  c3  c4  0  c4  c3
l
equation (6) becomes
ny  ny
nx   nx ny
u ( x, y )  bn sin e l   2bn sin
e l
sinh [c2 c3  bn ] (7)
l  
 l l
1
putting y=a, we have bn 
na
2 sinh
l
 nh 
sinh
nx  l  which is required solution.
Hence (7) reduces to u ( x, y )  sin  
l  sinh na 
 l 

Session 26
1. Obtain all possible solutions of two dimensional Laplace equation under steady state
conditions by method of separation of variables.
2. A infinitely long uniform plate is bounded by two parallel edges and an end at right angles
to them. The breadth is π. This end is maintained at a temperature u 0 at all points and other
edges are at zero temperature. Determine the temperature at any point of the plate in the
steady state.

Home Work

1. Find the steady-state temperature u(x , y) in a square plate 1 m on a side where u(x,
1)= x-x2 for 0 ≤ x ≤ 1 and u (x , y)=0 on the other three sides.

CO-3 Probability and Random Variables

Random Experiment: An experiment is called a Random Experiment if, when repeated

under the same conditions, it is such that the outcome cannot be predicated with certainty but
all possible outcomes can be determined prior to the performance of the experiment.
For Example: 1. Throwing of a die,
2. Tossing of a coin,
3. Drawing two playing cards from a pack of cards.
Sample Space: The set of all possible outcomes of a Random experiment is called the Sample
space and is represented by the symbol S.
For Example 1. When a coin is tossed the sample space is S  H , T  .

77
 2. When a six faced die is rolled the sample space is S  1,2,3,4,5,6.
Event: An event is subset of a sample space.
For Example: When a six faced die is rolled A  2,4,6 is a event and represent the
occurrence of an even numbers of dots.
Events are denoted by A, B , C ,  or E1 , E 2 , E3 , 
An event may be a subset that includes the entire sample space S called entire event, or a subset
of S called the null set and denoted by the symbol  , which contains no elements at all called
null event.
For instance, if we let A be the event of detecting a Microscopic organism by the naked eye in
a Biological experiment, then A   .
Also, if B   x / x is an even factor of 15 , then B must be the null set.
Complement of an Event: The complement of an event A with respect to S is the sub set
of all elements of S which are not in A . We denote the complement of A by the symbol A1

or A c or A .
Example: Let A be an event that an even number of dots occurred when a die is rolled then A1
is an event that an odd number of dots occurred.
Intersection of Two Events: the intersection of two events A and B denoted by the symbol
A  B , is the event containing all elements that are common to A and B.
Example: Let C be the event that a student selected at random is a second year student and M
be the event that student is a boy then C  M is the event of all second year boys.
Mutually Exclusive Events: Two events A and B are mutually exclusive, or disjoint if
A  B   , i.e., if A and B have no elements in common.
Example: In the die tossing experiment, if A  1,3,5 and B  2,4.6 then the events A and B
are mutually exclusive.
Union of Two Events: The union of two events A and B, denoted by the symbol A  B , is the
event containing all the elements that belong to A or B or both.
Example: In die tossing experiment, if A  3,6 and B  2,4,6 then A  B  2,3,4,6 and
it represent the event of getting an even number or a multiple of 3 dots. The following results
can be observed:
1) A     .
2) A    A.
3) A  A1   .
4) A  A1  S .
5) S 1   .
6)  1  S .
7) A1   A.
1

8) ( A  B)1  A1  B 1 .
9) ( A  B)1  A1  B 1 .
Classical Definition of Probability:

78
 If there are n outcomes mutually exclusive and equally likely outcomes of a random
experiment, out of which, ' s' outcomes are favorable for a particular E , then we define the
s Number of favourable outcomes of the exp eriment
probability of E , as P ( E )   .
n Number of total outcomes of the exp eriment
This probability is also know as probability of success of E .
In this experiment ' s' results are favorable to E, and hence the remaining 𝑛 − 𝑠 results are not
favorable to the event E . This set of unfavorable events denoted by E 1 or E c or E.
ns s
Probability of P( E c )   1   1  P( E ).
n n

Statistical Definition of Probability (or Relative Frequency Interpretation of

Probability):
Let m be the frequency of occurrence of the event A associated with the n independent trails
of the random experiment. Then probability of event A , denoted by the symbol P(A) ) is given
by
m
P( A)  Lt .
n  n

m
We may note that is the relative frequency of the event A in 𝑛 trails. If 𝑛 is very large then
n
m
the relative frequency is very close to actual probability.
n
Axiomatic Definition of Probability:
Probability is a number that is assigned to each member of a collection of events from a random
experiment that satisfies the following properties.
If S is the sample space and E is any event in a random experiment,
(i) 0  P ( E )  1 for each event e in S.
(ii) P ( S )  1.
(iii) If E1 and E2 are any mutually exclusive events in S, then P( E1  E 2 )  P( E1 )  P( E 2 ).
Example 1: A box contains 25 parts of which 10 are defective. Two parts are being drawn
simultaneously in a random manner from the box. The probability of both parts being good is
[2014 Mech]
A) 7/20 B) 42/125 C) 25/29 D) 5/9
Solution: Number of ways of drawing 2 parts from 25 parts=25c2
Number of ways of drawing 2 good parts from the 15 good parts=15c 2
Probability that both parts are good= =
Example 2: In a housing society, half of the families have a single child per family, while the
remain half have two children per family. The probability that a child picked at random, has a
sibling [2014: 1 mark set -1 E.C.E ]
Solution:
The child picked at random will have a sibling if the family has two children
Probability of this event=1/2.
Example 3: An unbiased coin is tossed an infinite number of times. The probability that the
fourth head appears at the 10th toss is [2014: 1 mark set -3 E.C.E ]
A) 0.067 B) 0.073 C) 0.082 D) 0.091

79
Solution:
Total number of possibilities for the first ten slips is 2 10=1024.
For the fourth head to occur on 10th slip.
We need first 3 heads to occur in the first 9 slips. This is given by
9C3=84.
There is only one way for 4th head occur on 10th slip =84*(1/1024)=21/256=0.082.
(or) For the 4th head to occur at the 10th toss, you have to first get 3 heads and 6 tails in the
1st 9 toss, and then a head at the 10th toss.
So prob. = (9C3)(.5)^3 (1-.5)^6 (0.5)=.082.
Example: 4 A fair dice is tossed 10 times. What is probability that nly the frst two tosses will
yield heads
A) (1/2)^2 B) 10C2 (1/2)^2 C) (1/2)^10 D) 10C2(1/2)^10 [2009: 1 mark
ECE]
Solution: (1/2)2(1/28=(1/2)10

Practise Problems
1) A pair of fair dice is tossed. Find the probability of getting
(a) a total of 8;
(b) at most a total of 5.
(c) same number of dots on both dice
(d) at least a total of 10
2) Interest centers around the life of an electronic component. Suppose it is know that the
probability that the component survives for more than 6000 hours is 0.42. Suppose also that
the probability that the component survives no longer than 4000 hours is 0.04.
a) What is the probability that the life of the component is less than or equal to 6000 hours
b) What is the probability that the life of the component is greater than 4000 hours
3) In a poker hand consisting of 5 cards, find the probability of holding
(a) 3 aces
(b) 4 hearts and 1 club
(c) cards of same suit
(d) 2 aces and 3 jacks

Addition Theorem of Probability:

If A and B are two events, then
PA  B )  P ( A)  P ( B )  P ( A  B ).
Note:
1) If A and B are mutually exclusive then PA  B )  P ( A)  P ( B ).
2) For three events A, B and C then
P ( A  B  c )  P ( A)  P ( B )  P (C )  P ( A  B ).  P ( A  C )  P ( B  C )  P ( A  B  C ).
3) If A, B and C are mutually exclusive the PA  B  C )  P ( A)  P ( B )  P (C ).
4) If A1 , A2 , , An are n mutually exclusive events then
P( A1  A2    An )  P( A1 )  P( A2 )    P( An ).
Example: Suppose that in a senior college class of 500 students it is found that 210 smoke,
258 drink alcoholic beverages, 216 eat between meals, 122 smoke and drink alcoholic

80
beverages, 83 eat between meals and drink alcoholic beverages, 97 smoke and eat between
meals, and 52 engage in all three of these bad health practices. If a member of this senior class
is selected at random, find the probability that the student
a) Smokes but does not drink alcoholic beverages.
b) eats between meals and drinks alcoholic beverages but does not smoke;
c) Neither smokes nor eats between meals.
d) Probability that the student does not have any habit
Solution: Let A, B, and C be the events that the student selected at random is found to be
smoke, drink, alcoholic beverages and eat between meals, respectively. From the given data
P(A)=210/500, P(B)=258/500, P(C0=216/500, P(A∩ 𝐵)=122/500,
P(B∩C)=83/500, P(A∩C)=97/500 and P(A∩ B ∩ C)=52/500
a) Probability that the student selected at random smoke but does not drink alcoholic
beverages
=P(A∩ B)=P(A)- P(A∩ 𝐵) =21/500-(122/500)=88/500
b) Probability that the student selected at random eat between meals and drink alcoholic
beverages but does not smoke

=P(C∩ 𝐵 ∩ 𝐴̅)=P(𝐴̅ ∩ 𝐵 ∩ 𝐶)=P(B∩ 𝐶)-P(A∩B∩C)=(83/500)-(52/500)=31/500.

c) Probability that the student neither smokes nor eats between meals
=P(A ∩ C)=P(A ∪ C)=1-P(A∪ 𝐶)=1-[P(A)+P(C)-P(A∩ 𝐶)]=1-[(210/500)+(216/500)-
(97/500)]=171/500.
d) Probability that the student does not have nay habit
= P(A ∩ B ∩ C)=P(𝐴 ∪ 𝐵 ∪ 𝐶)=1-P(A∪ 𝐵 ∪ 𝐶)=1-[P(A)+P(B)+P(C)-P(A∩B)-P(B∩C)-
P(C∩A)+P(A∩ 𝐵 ∩ 𝐶)]=1-(434/500)=66/500.

Practise Problems

1. A final year student after being interviewed at two companies, he assesses that his
probability of getting an offer from company A is 0.8 and the probability that he gets offer
from company B is 0.6. If, on the other hand he believes that the probability that he will get
offers from both companies is 0.5. Obtain the probability that he will get
(i) at least one offer from these two companies
(ii) offer from neither company
(iii) offer from company A only
(iv) offer from only one company

2. Suppose that in the maintenance of a large medical records file for insurance purposes the
probability of an error in processing is 0.0010, the probability of an error in filing is 0.0009,
the probability of an error in retrieving is 0.0012, the probability of an error in processing
as well as filing is 0.0002, the probability of an error in the probability of an error in
processing as well as retrieving is 0.0003, and the probability of an error in processing and
filing as well as retrieving is 0.0001. a) Obtain the probability of making at least one of these
errors?
b) Obtain the probability of making none of these errors?
3. From the past experience a stockbroker believes that under present economic conditions a
customer will invest in tax-free bonds with a probability of 0.6 will invest mutual funds with

81
probability of 0.3 and will invest in both tax-free bonds and mutual funds with probability of
0.15. At this time, find the probability that a customer will invest
a) In either tax-free bonds or mutual funds
b) in neither tax-free bounds nor mutual funds
c) in only one investment
4. In a high school graduating class of 100 students, 54 studied mathematics, 69 studied history,
and 35 studied both Mathematics and history. If one of these students is selected at random,
find the probability that
a) the student took Mathematics or History.
b) the student did not take either of these subjects.
c) the student took History but not Mathematics.
d) the student took only one of History or Mathematics.
Conditional Probability:
The Conditional Probability of B, given A, denoted by P(B/A) is defined by
( ∩ )
P(B/A) = , provided P(A) > 0.
( )
Independent Events
If the occurrence of B had no impact on the odds of occurrence of A, then A and B are said to
be independent.

Example: If P(A)=0.65, P(B)=0.40 and P(C∩D)=0.24, are the events C and D independent?
Two events A and B are independent if and only if P ( B / A)  P ( A) or P ( A / B )  P ( B ),
provided the existence of the conditional probability.

Example 1: The chance of a student passing an exam is 20%The chance of passing an exam
and getting above 90% marks in it is 5% given that the student passes the examination, the
probability that the student gets above 90% marks is [2015 Mech set 2]
A ) 1/18 B) 1/4 C) 2/9 D) 5/18

Solution: Let A and B denote the events of a student passing an exam and a student getting
above 90% marks in the exam respectively.
P(A)=20/100=0.2, P(B)=5/100=0.05.
Given that a student passes the examination, the probability that the students gets above 90%
marks
=P(B/A)=0.05/0.2=1/4.

Example 2: P(X) = 1/4, P(Y) = 1/3, P(X  Y ) = 1/12 The value of P(Y/X) is [2015 Mech
set 3]
A) 1/4 B) 4/25 C) 1/3 D) 29/50
1
𝑌 𝑃(𝑌 ∩ 𝑋) 12
𝑃 = = = 1/3.
𝑋 𝑃(𝑋) 1
4

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Multiplicative Rule
If in a experiment the events A and B can both occur, then P ( A  B )  P ( B / A) P ( A) provided
P(A) > 0. We can also write P ( A  B )  P ( A / B ) P ( B )
In other words, it does not matter which event is referred to as A and which event is referred
to as B.

Note:
1) Two events A and B are independent if and only if P ( A  B )  P( A) P ( B ) .
eg: Let A be the event that raw material is available when needed and B be the event that
the matching time is less than one had. If P(A)=0.8 and P(B)=0.7. What is P(A∩B).
2) If A, B, C are any three events then the multiplicative rule
P ( A  B  C )  P ( B / C ) P ( A) P (C / A  B ) .
3) If A, B, and C are independent events if and only if P ( A  B  C )  P ( A) P ( B ) P (C ).

Example1: Two cards are drawn at random from an ordinary deck of 52 playing cards. What
is the probability of getting two aces if
a) The first card is replaced before the second card is drawn;
b) The first card is not replaced before the second card is drawn?
Solution:
a) Since there are four aces among the 52 cards, we get (4/52).(4/52)=1/169.
b) Since there are only three aces among the 51 cards that remain after one ace has been
removed from the deck, we get
4 3 1
. =
52 51 221
Note that . ≠
So independence is violated when the sampling is without replacement.

Example2:.Three vendors are asked to supply a very high precision component. The respective
probabilities of their meeting the strick design specifications are 0.8, 0.7 and 0.5 Each vendor
supplies one component .The probability out of total three components supplied by the vendors
atleast one will meet the design specifications is --------- [ 2015 Mech set 2 ]
Let A, B, C be the event that the high precision component supplied by the three vendors
meets the design specifications. The events A, B and C are independent.
Given that P(A)=0.8, P(B)=0.7 and P(C)=0.5
Probability that at least one meet the design specifications=1-Probability that none of them
meet the design specifications
=1-P(A ∩ B ∩ C) = 1 − P(A)P(B)P(C)=1-(0.2)(0.3) (0.5)=0.97.

Example: 3 . A coin is tossed thrice .Let X be the event that head occurs in each of the first
two tosses. Let Y be the event that tail occurs on the third toss. Let Z be the event that two
tails occur in three tosses .Based on the information which out of the following statements is
true [2015 Mech set 3]
A) X and Y are not independent
B) Y and Z are dependent

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C) Y and Z are independent
D) X and Z are independent
Ans: C
Example: 4 A fair coin is tossed till a head appears in the first time.The probability that the
number of required tosses is odd, is [ 2012: 2 marks ECE]
A) 1/3 B) 1/2 C) 2/3 D) ¾
Add upto probability of the coin coming up head for the first time 1, 3,5, …
P0=(1/2)1+(1/2)3+(1/2)5+….=(1/2)/(1-(1/2)2)=(1/2)/(1-1/4)=2/3.
Example: 5 A fair dice is tossed till two times. The probability that the second toss results in
a value that is higher than the first toss is [2011: 2 marks ECE]
A) 2/36 B) 2/6 C) 5/12 D) ½

Solution:
(1/6)*(5/6)+(1/6)*(4/6)+(1/6)*(3/6)+(1/6)*(2/6)+(1/6)*(1/6)=15/36=5/12 (or)
Pr(Second >first)+Pr(Second<first)+Pr(Second=first)=1
By symmetry
P(Second>first)=1-P(Second=first)
P(Second>First)=(1-1/6)/2=5/12.

Example: 6 A fair coin is tossed independently four times. The probability of the event ‘the
no of times heads show up is more than the no of times tails show up’ is [2010: 2 marks ECE]

A) 1/16 B) 1/8 C) 1/4 D) 5/16

Solution:

1 tail, 3Heads

4c3(1/2)4+(1/2)4=(4/16)+(1/16)=5/16.

Example: 7 A 1-h rain fall of 10cm magnitude at a station has a return period of 50
years.The probability that a 1-h rain fall of 10cm magnitude or more will occur in each of two
successive years is

A) 0.004 B) 1.0 C)1.5 D) 2.0

Solution: (1/50)*(1/50)=1/2500=0.004

Practise Problems
1) The Probability that a regularly scheduled flight departs on time is P(D)=0.83; the
probability that it arrives on time is P(A)=0.82; and the probability that it departs and arrives
on time is P ( D  A)  0.78. Find the probability that a plane
(a) Arrives on time given that it departed on time.
(b) Departed on time given that it has arrived on time,
(c) Neither departed on time nor arrived on time.

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2) A small town has one fire engine and one ambulance available for emergencies. The
probability that the fire engine is available when needed is 0.98, and the probability that the
ambulance is available when called is 0.92. In the event of an injury resulting from a burning
building, find the probability that
a) Both the ambulance and the fire engine will be available.
b) Ambulance or fire engine available.

3) In the senior year of a high school graduating class of 100 students, 42 studied mathematics,
68 studied psychology, 54 studied history, 22 studied both mathematics and history, 25 studied
both mathematics and psychology, 7 studied history but neither mathematics nor psychology,
10 studied all three subjects and 8 did not take any of the three. If a student is selected at
random, find the probability that
a) a person is taking at least one subject
b) a person enrolled in psychology takes all three subjects.
c) a person not taking psychology is taking both history and mathematics.

4) A manufacturer of a flu vaccine is concerned about the quality of its flu serum. Batches of
serum are processed by three different departments having rejection rates of 0.10, 0.08, and
0.12 respectively. The inspections by the three departments are sequential and independent.
a) what is the probability that a batch of serum survives the first departmental inspection but is
rejected by the second department?
b) what is the probability that a batch of serum is rejected by the third department?
c) what is the probability that a batch of serum survives all three departmental inspection?
Session 27
Baye ‘s Theorem
Total Probability:
If the events B1 , B2 , , B K constitute a partition of the sample space S such that
P( Bi )  0 for i  1,2, , k , then for any event A of S,
k k
P ( A)   P ( Bi  A)   P ( Bi ) P ( A / Bi ).
i 1 i 1

Example: 1 A group contains equal no of men and women of those 20 % of the men, 50 % of
women are unemployed. If a person is selected at random from these. The probability of
selected person being employed is ----- [2014 Mech]
Solution:
Let E be the event that the probability of selected person being employed is

85
If man is selected then probability of him to be employed=20/100.
Similarly if woman was selected the probability of her to be employed=50/100
So
probability=P(M).P(E/M)+P(W).P(E/W)=(1/2)*(20/100)+(1/2)*(50/100)=(1/10)+(5/20)=7/2
0.
Required probability =1-P(E)=1-0.35=0.65.
Baye’s Rule:
If the events B1 , B2 ,, BK constitute a partition of the sample space S such that
P( Bi )  0 for i  1,2, , k , then for any event A in S, such that P(A)≠0,
P ( Br  A) P ( Br )
P ( Br / A)  K
 k
, for  1, 2,..., k
 P( B  A)  P( B ) P( A / B )
i 1
i
i 1
i i

Example l: The probability that student knows the correct answer to a multiple choice
question is 2/3. If the student does not know the answer, the student guesses the answer. The
probability of guessed answer is correct is ¼ . Given that student has answered the question
correctly. The conditional probability that student knows the correct answer is [2013 mech]
A) 2/3 B) 3/4 C) 5/6 D) 8/9
Solution: Let A be the event that the student knows the correct answer. Then A represent the
event that the student guesses.
P(A)=2/3, P(A)=1-P(A)=1/3
Let E be the event that the answer is correct
Given that P(E/A)=1/4
Since the student answers correctly when he knows the correct answer P(E/A)=1
Probability that the student knows the correct answer given that he answers correctly
P(A ∩ E) P(A ∩ E)
P(A/E) = =
P(E) P(A ∩ E) + P(A ∩ E)
E
P(A)P( )
= A = 8/9
E E
P(A)P( ) + P(A)P
A A
Problems to be discussed by the Faculty:

1. A box contains 5 red marbles and 4black marbles. Another box contains 3 red marbles
and 6 black marbles. A box is selected at random and a marble is drawn from it.
(a) What is the probability that the marble drawn is red?

86
 (b) If the marble drawn is known to be red what is the probability it came from first
box?
2. At an electronic plant, it is known from past experience that the probability is 0.83 that
a new worker who has attended the company’s training program will meet the
production quota and that the corresponding probability is 0.35 for a new worker who
has not attended the company’s training program. If 80% of all new workers attend the
training program, what is the probability that a new worker will meet the production
quota? If a new worker meets the production quota, what is the probability that he
attended the training programme?

3. The members of a consulting firm rent cars from three rental agencies: 60% from
Agency 1, 30 % from Agency 2 and 10% from Agency 3. If 9% of the cars from
Agency 1 need a tune-up, 20 % of the cars from agency 2 need a tune up and 6% of
the cars from agency 3 need a tune-up,
(a) What is the probability that a rental car delivered to the firm will need a tune -up?
(b) If the car delivered need a turn-up, what is the probability that it was from
(i) agency1 (ii) agency2 (iii) agency3
Home Work
1. In a certain state, 25% of all cars emit excessive amounts of pollutants. If the
probability is 0.99 that a car emitting excessive amounts of pollutants will fail
the state’s emission test, and the probability is 0.17 that a car not emitting
excessive amounts of pollutants will nevertheless fail the test,
(a) What is the probability that a car fails the emission test.
(b) What is the probability that a car that fails the emission test actually emits
excessive amounts of pollutants?
2. A box contains 4 green marbles and 2 blue marbles and another box contains 3
green marbles and 3 blue marbles. A die is tossed and if the outcome is 1 or 2
first box is selected, otherwise second box is selected and a marble is drawn
from it.
(a) What is the probability that the marble drawn is green?
(b) If the marble drawn is blue, what is the probability that it come from second box
3. The probability that a car accident is due to faulty brakes is 0.04, the
probability that a car accident is correctly attributed to faulty brakes is 0.82,
and the probability that a car accident is incorrectly attributed to faulty brakes
is 0.03. What is the probability that :

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 a) A car accident will be attributed to faulty brakes;
b) A car accident attributed to faulty brakes was actually due to faulty brakes .

Session 28

RANDOM VARIABLE:
A random variable is a function that associates a real number with each element in the sample
space.

Ex: The testing of three of electronic components for defectives is a random experiment. The
associated sample space is
S  NNN , NND , NDN , DNN , NDD , DND , DDN , DDD  where N denotes non defective
and D denotes “defective”. If X is the number defectives. Then for each point in the sample
space, X associates real numbers 0, 1, 2, or 3.

Types of Random Variables: Random variables are of two types: Discrete random variable
and continuous random variable.

Discrete Random Variable: A random variable is called a discrete random variable if its set
of possible outcomes is countable. If X is the number of defective components out of 3
components. Then X is a discrete random variable.

Discrete probability distribution function: A discrete random variable assumes each of its
values with a certain probability. The table showing the values of the random variable and
associated probabilities is called the discrete probability distribution.
Ex: In the case of tossing a coin three times, the variable X, representing the number of heads
has the following probability distribution.
x 0 1 2 3
f(x)=P(X=x) 1/8 3/8 3/8 1/8

Probability function: The set of ordered pairs (x, f(x)) is called the probability function or
probability mass function or probability distribution of the discrete random variable X, if for
each possible outcome x,
1. f ( x )  0.
2.  f ( x)  1
x

3. P ( X  x)  f ( x).
Cumulative distribution: The cumulative distribution function F(x) of a discrete random
variable X with probability distribution f(x) is
F ( x )  P ( X  x )   f (t ) for    x .
tx

Properties:
(i) F (  )  0
(ii) (ii) F ( )  1
(iii) (iii) P ( a  X  b )  F (b)  F ( a )

88
 (iv) F(x) is a non decreasing function.

Example:
Given that f(x)=k/2x is a probability distribution for a random variable that can take on the
values x=0, 1, 2, 3 and 4. Then, (a) Find K, (b) Find the Cumulative probability distribution
F(x)
Solution: (a) since f(x) is the probability mass function

f(x) = 1

 f(0)+f(1)+f(2)+f(3)+f(4)=1
 k/1+k/2+k/4+k/8+j/16=1
 K[31/16]=1
 K=16/31.
(b) F(0)=f(0)=k=16/31
F(1)=f(0)*f910=K+(K/2)=(3/2)K=24/31
F(2)=f(0)+f(1)+f(2)=K+(K/2)+(K/4)=(7/4)*K=28/31
F(3)=f(0)+f(1)+f(2)+f(3)=K+(K/2)+(K/4)+(K/8)=30/31.
F(4)=f(0)+f(1)+f(2)+f(3)+f(4)=K+(K/2)+(K/4)+(K/8)+(K/16)=1
c) f(2)=F(2)-F(1)=(28/31)-(24/31)=4/31.

Continuous Random Variable: When a random variable can take on values on a continuous
scale, it is called a continuous random variable.

Ex: Le X be the random variable defined by the waiting time, in hours, between successive
speeders spotted by a radar unit. The random variable X take on values x for which x  0.

Continuous Probability Distribution: A continuous random variable has a probability of

zero of assuming exactly any of its values. Consequently, its probability distribution cannot be
given in tabular form.

Density Function: In dealing with continuous variables, f(x) is usually called the probability
density function or simply the density function of X.
The function f(x) is a probability density function for the continuous random variable X,
defined over the set of real numbers R, if
1) f ( x )  0, for all x  R

2)  f ( x)dx  1

b
3) P ( a  x  b)   f ( x )dx.
a

Cumulative Distribution Function (CDF): The cumulative distribution function F(x) of a

continuous random variable X with density function f(x) is
x
F ( x)  P( X  x)   f (t )dt ,

for    x  

89
 dF ( x )
Note: P ( a  X  b)  F (b)  F ( a ) and f ( x )  , if the derivate exits.
dx
Example:
An important factor in solid missile fuel is the particle size distribution. Significant problems
occur if the particle sizes are too large. From the production data in the part, it has been
determined that the particle size (in micrometers) distribution is characterized by
f(x) =3x -4, x>1
= 0, elsewhere
a) Verify that this is a valid density function
b) Evaluate the Cumulative distribution function F(x)
c) What is the probability that a random particle from the manufactured fuel exceeds 4
micrometers?
d) What is the probability that a random particles’ size is between 2 and 4 micrometers?
Solution:
a) f(x) is a valid density function if ∫ 𝑓(𝑥)𝑑𝑥 = 0, x≤X
𝑥
𝑓(𝑥)𝑑𝑥 = 3𝑥 =3 =1
−3

Therefore, f(x) is a valid density function.

b) Cumulative distribution function F(x) is given by

f(x)=P(X≤ 𝑥)= ∫ 3x =1−
0, x < 1
∴ F(x) =
1 − ,x ≥ 1
c) Probability that a random particle’s size exceeds 4 micrometers
=P(X>4) = 1−F(4) = 1−(1−(1/43)=0.0156.
d) Probability that a random particle’s size is between 2 and 4 micrometers
=P(2≤X≤4)=F(4)-F(2)=(1-1/4 3)-(1-1/23)=1/8-1/64 =0.109.

Problems to be discussed by the faculty:

1. For each of the following, determine whether the given values can serve as the values
of a probability distribution of a random variable with the range x=1, 2, 3 and 4
(a) f(1) = 0.25, f(2) = 0.75, f(3) = 0.24 and f(4) = -0.25
(b) f(1) = 1/19, f(2) = 10/19, f(3) = 2/19 and f(4) = 4/19
2. Find the value of c so that the following function can serve as the probability
distribution of a random variable
(a) f(x) = c , x = 1, 2, 3,.........

3. A shipment of 8 similar microcomputers to a retail outlet contains 3 defectives. If a

school makes a random purchase of 2 of these computers, find the probability distribution
for the number of defectives. Also determine the mean and variance of X.
4. Consider the density function

90
  k x , 0  x  1
f ( x)  
 0, elsewhere

a) Evaluate k.
b) Evaluate P(0.3<X<0.6) using the density function
c) Determine mean and variance of X

5. The shelf life ( in hour) of a certain perishable packaged food is a random variable whose
,
, 𝑓𝑜𝑟 𝑥 > 0
probability density function is given by 𝑓(𝑥) = ( ) .
0, 𝑒𝑙𝑠𝑒 𝑤ℎ𝑒𝑟𝑒
Find the probabilities that one of these packages will have a shelf life of
(a) at least 200 hours, (b) at most 100 hours, (c) anywhere from 80 to 120 hours.
Home work

1. An urn contains four balls numbered 1, 2, 3 and 4. If two balls are drawn from the urn
at random and Z is the sum of the numbers on the two balls drawn, find the probability
distribution of Z.
𝑘𝑒 , 𝑓𝑜𝑟 𝑥 > 0
2. If X has the probability density 𝑓(𝑥) = 0, 𝑒𝑙𝑠𝑒 𝑤ℎ𝑒𝑟𝑒 . Find (i) k, (ii) P (0.5 ≤

𝑋 ≤ 1) (iii) P(X < 1) and (iv) 𝑃(𝑋 ≥ 1)

3. The probability density function of the random variable X is given by 𝑓(𝑥) =
, 𝑓𝑜𝑟 0 < 𝑥 < 4
√
0, 𝑒𝑙𝑠𝑒 𝑤ℎ𝑒𝑟𝑒 . Find (i) the value of c (ii) P(X<1) and (iii) 𝑃(𝑋 ≥ 1).

Session 29

BINOMIAL DISTRIBUTION AND POISSON DISTRIBUTION

Bernoulli process: An experiment often consists of repeated trials, each with two possible
outcomes that may be labelled success or failure. As an example, the testing of items as they
come off an assembly line, where each test or trial may indicate a defective or non defective
item. We may choose to define either outcome as a success. The process is referred to as a
Bernoulli process. Each trial is called a Bernoulli trial.

Properties of Bernoulli Process:

The Bernoulli process must possess the following properties:

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1. The experiment consists of n-repeated trials.
2. Each trial results in an outcome that may be classified as a success or a failure.
3. The probability of success, denoted by p, remains constant from trial to trial.
4. The repeated trials are independent.
Consider a sequence of 𝑛 independent trials of a Bernoulli process. Let X be the number of
success in n-Bernoulli trials and it is a random variable. The distribution of the random variable
X is called Binomial distribution.

Definition: A Bernoulli trial can result in a success with probability p and a failure with
probability q =1−p. Then the probability distribution of the Binomial random variable X, the
number of success in n-independent trials, is

b(x; n, p) =   p x q n  x ,
n
x = 0, 1, , n, p + q = 1.
 x

Note
1) We write, X ~ b(n, p) to denote that X follows binomial distribution with parameters n and
p
2) The mean of the binomial distribution is ‘np’
3) The variance of the binomial distribution is ‘npq’
4) The standard deviation of binomial distribution is 𝑛𝑝𝑞
5) In binomial distribution, the mean is always greater than the variance.
Example: It has been claimed that in 60% of all solar-heat installations the utility bill is
reduced by at least one-third. Accordingly, what are the probabilities that the utility bill will
be reduced by at least one-third in
a) Four of five installations
b) at least four of five installations
c) at the most two installations
d) what are the mean and variance of the number of installations
Solution: Let X be the number of solar installation where the utility bill is reduced by at least
one third. Then the distribution of X is binomial with n=5 and P=0.6
P(X=x) =5 (0.6)x(0.4)5-x, x=0,1,2,3,4,5
a) Probability that in four of 5 installations utility bill is reduced by one third is
=P(X=4)= 5 (0.6)4(0.4)5-4=0.259
b) Probability that in at least 4 installations utility bill is reduced by one third is

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 =P(X≥4)=P(X=4)+P(X=5)= 5 (0.6)4(0.4)5-4+5 (0.6)5(0.4)5-5
c) Probability that in at most two installations utility bill is reduced by one third is
=P(X≤2)=P(X=0)+P(X=1)+P(X=2)
=5 (0.6)0(0.4)5-0+5 (0.6)1(0.4)5-1+5 (0.6)2(0.4)5-2
d) Mean=np=5(0.6)=3
Variance=np(1-p)=5(0.6)(0.4)=1.2.
POISSON DISTRIBUTION
When n is large and p is small, binomial probabilities are often approximated by means of the
Poisson distribution with the parameter λ equal to the product np i.e., Poisson distribution is
used in case of rare events.
Experiments yielding numerical values of a random variable X, the number of outcomes
occurring during a given time interval or in a specified region (Here n is not known and
information regarding the number of occurrences of event is known), are called Poisson
experiments.
The number X of outcomes occurring during a Poisson experiment is called a Poisson random
variable and its probability distribution is called the Poisson distribution.

Definition: The probability distribution of the Poisson random variable X, representing the
number of outcomes occurring in a given time interval or specified region denoted by t, is
e (λt)
P(x; λt) = , x = 0,1,2, … ..
x!
where λ is the average number of outcomes per unit time, distance, area or volum and
e=2.71828

Note: Both the mean and variance of the Poisson distribution P(x; λt) are λt.
1) If the time t is one unit then
e (λ)
P(x; λ) = , x = 0,1,2, … ..
x!
2) We write, X ~ P(x, ) to denote that X follows Poisson distribution with parameter .
The following are the some of the examples of random variables following Poisson
distribution:
 The number of customers arrived during a time period of length t.
 The number of telephone calls per hour received by an office.
 The number of typing errors per page.

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  The number of accidents occurred at a junction per day.
Example 1: If a bank received on the average 6 bad checks per day, what are the
probabilities that it will receive
a) 4 bad checks on any given day?
b) 10 bad checks over any 2 consecutive days
c) No bad check on any given day
d) What are the mean and variance of the number of bad check per day?
Solution: Let X be the number of bad checks received per day. Then the distribution of X is
Poisson with parameter λ=6.
e 6
P(X = x) = , x = 0,1,2, …
x!

a) P(4 bad checks on any given day) = P(X = 4) = = 0.134

!

. b) P(10 bad checks over any 2 consecutive days) = P(X = 2) = , x = 0,1,2, ….

!

(Here λ = 12)

. c) P(no bad check on any day) = P(X = 0) = =e

!

d) Mean and variance of the number of bad checks per day = λ = 6.

Example 2: In the inspection of tin plate produced by a continuous electrolytic process, 0.2
imperfections is spotted per minute, on average. Find the probabilities of spotting
a) One imperfection in 3 minutes
b) at least two imperfections in 5 minutes
c) at most one imperfection in 15 minutes
Solution: Let X be the number of imperfections spotted per minute. Then the distribution of
X is Poisson with λ=0.2
a) Average number of imperfections in 3 minutes=λt=3(0.2)=0.6
( . )
Probability of one imperfection in 3 minutes= = 0.329
!

b) Average number of imperfections in 5 minutes= λt = 5(0.2) = 1.0

Probability of at least two imperfections per 5 minutes
=P(X≥2)=1-P(X<2)=1-P(X=0)-P(X=1)=1-e-1-(e-1(1.0)/1)=0.264.

c) Average number of imperfections in 15 minutes= λt = 15(0.2) = 3.0

Probability of at most one imperfection in 15 minutes

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 =P(X≤1)=P(X=0)+P(X=1)=e-3+((e-331/1!)=0.199.
Problems to be discussed by the faculty:
1. The mean and variance of a binomial variable are 5 and 5/2 respectively, find the
values of the parameters n and p.
2. An automobile engineer claims that 1 in 10 automobile accidents are due to driver
fatigue. If a random sample of 5 accidents observed, what is the probability that
(a) 3 accidents are due to drive fatigue, (b) at most 1 due to drive fatigue, (c) at least 4
accidents are due to driver fatigue.
3. In a Poisson distribution with parameter 𝜆, P(X=1) = P(X=2), find the mean and
variance of the distribution.
4. If 2% of the books found at a certain bindery have defective findings. Out of 400
books found by this bindery, what is the probability that
(a) 3 books have defective bindery
(b) at least one book has defective bindery
(c) at most one book has defective bindery
5. The average number of trucks arriving on any one day at a truck depot in a certain city is
known to be 12. What is the probability that on a given day fewer than 2 trucks will arrive at
this depot.

HOME WORK
1) In testing a certain kind of truck tire over a rugged terrain, it is found that 25% of the
trucks fail to complete the test run without a blowout. Of the next 15 trucks tested, what
is the probability that
(a) 3 have blowouts
(b) fewer than 2 have blowouts
(c) more than 14 blowouts.
2) A traffic control engineer reports that 75% of the vehicles passing through a checkpoint
are from within the state. What is the probability that fewer than 4 of the next 9 vehicles
are from out of state?
3. In a certain industrial Facility accident occur infrequently. It is known that the probability
of an accident on any given day is 0.005 and accidents are independent of each other.
a) What is the probability that in any given period of 400 days there will be an accident on
one day.
b) What is the probability that there are at most three days with an accident?
4. In a manufacturing process where glass products are produced, defects or bubbles occur,
occasionally rendering the piece undesirable for marketing. It is known that, on average, 1 in

95
every 1000 of these items produced has one or more bubbles. What is the probability that a
random sample of 8000 will yield
a) Fewer than 7 items possessing bubbles
b) None will possess a bubble
Session 30 and 31
Normal Distribution – Its properties and importance
The most important continuous probability distribution in the entire field of statistics is the
normal distribution. It s graph, called the normal curve, is the bell shaped curve, which
describes approximately many phenomena that occur in nature, industry, and research. Physical
measurements in areas such as meteorological experiments, rainfall studies, and measurements
of manufactured parts are often more than adequately explained with a normal distribution.

The basic from of normal distribution is that of a bell, it has single mode and is symmetric
about its central values. The flexibility of using normal distribution is due to the fact that the
curve may be centered over any number on the real line and it may be flat or peaked to
correspond to the amount of dispersion in the values of random variable.
The Normal Distribution (N.D.) was first discovered by De-Moivre as the limiting form of the
binomial model in 1733. The normal distribution is often referred to as the Gaussian
distribution, in honour of Karl Friedrich Gauss (1777-1855), who also derived its equation from
a study of errors in repeated measurements of the same quantity.
Normal distribution is often used to model the distribution of discrete random variable as well
as the distribution of other continuous random variables. Normal distribution provided the basis
for which much of the theory of inductive statistics is founded.

Definition

96
A random variable X is said to follow a Normal Distribution with parameter mean (µ) and
variance (𝜎 ) if its density function is given by the probability law
1 µ
f(x) = n(x; u, σ) = e , −∞<x<∞
σ√2π
Symbolically we can represent the distribution of normal variate as X ~ N(, 2)
The Properties of normal probability curve
1. The mode which is point on the horizontal axis where the curve is a maximum, occurs
at x=µ. Hence the mean, median and mode of normal distribution are equal.
2. The curve is symmetric about a vertical axis through the mean µ.
3. The curve has its points of inflexion at x=µ±𝜎, it concave downward if µ−𝜎<X< µ+𝜎,
and is concave upward otherwise.
4. The curve approaches the horizontal axis asymptotically as we proceed in either
direction away from the mean.
5. The total area under the curve and above the horizontal axis =1.
6. The graph of the normal distribution depends on two factors – the mean and the standard
deviation. The mean of the distribution determines the location of the center of the
graph, and the standard deviation determines the height and width of the graph. When
the standard deviation is large, the curve is short and wide; when the standard deviation
is small, the curve is tall and narrow. All normal distributions look like a symmetric,
bell-shaped curve, as shown below.

The curve on the left is shorter and wider than the curve on the right, because the curve
on the left has a bigger standard deviation.
7. Linear combination of independent normal variates is also a normal variate
8. The total area under the normal curve (∫ 𝑓(𝑥)𝑑𝑥 = 1) is distributed as follows

 (𝜇 − 𝜎) < 𝑥 < (𝜇 + 𝜎)covers 68.26% of the area

 (𝜇 − 2𝜎) < 𝑥 < (𝜇 + 2𝜎)covers 95.44% of the area
 (𝜇 − 3𝜎) < 𝑥 < (𝜇 + 3𝜎)covers 99.74% of the area, and it can be represented as
follows

97
Standard normal distribution:
The distribution of a random variable with mean ‘0’ and variance ‘1’ is called a standard normal
distribution. If Z is a standard normal variate then Z~N(0,1).
The probability density function of the standard normal variate Z is given by the probability
law
1
𝑓(𝑧) = 𝑒 , −∞ < 𝑧 < +∞
√2𝜋
The standard normal distribution, N (0, 1), is very important because probabilities of any
normal distribution can be calculated from the probabilities of the standard normal distribution.
Note:
X 
1. If X is a normal random variable with mean  and standard deviation , then Z 

x1   x2  
is a standard normal random variable and hence P( x1  X  x2 )  P( Z )
 
2. Suppose Z ~ N(0, 1) is standard normal variate then by using the standard normal
distribution area tables, we can calculate the various probabilities as explained below:
i) P(Z<a)P(Z≤a). This probability can be read from the table and is described in the
following figure

ii) P(Z>b). This probability can be represented by using the following graph of standard
normal distribution and it cannot be read directly from the standard normal tables

98
 P(Z>b)=1- P(Z≤b), where P(Z≤b) available directly from table.
iii) P(a≤Z≤b) . This probability can be represented by using the following graph of
normal distribution

P(a≤Z≤b) = P(Z≤b) – P(Z≤a), where P(Z≤b) and P(Z≤a) are available directly from
standard normal tables.
The Normal approximation to the Binomial distribution
Given X is a random variable which follows the binomial distribution with parameters
n and p, then the limiting form of the distribution is standard normal distribution i.e.,
𝑍= , 𝑤ℎ𝑒𝑟𝑒𝑞 = 1 − 𝑝 provided, if n is large and p is not close to 0 or 1 is a
√

standard normal variate.. If both np and nq are greater than 5, the approximation will
be good.

Example 1:
With an eye toward improving performance, industrial engineers studied the ability of scanners
to read the bar codes of various food and household products.The maximum reduction in
power, just before the scanner cannot read the bar code at a fixed dictionary is called the
maximum attenuation. This quantity, measured in decibles, varies from product to product:
After collecting the data, the engineers decided to model the variation in maximum attenuation
as a normal distribution with mean 10.1 dB and standard deviation 2.7 dB.
a) For the next food product, what is the probability that its maximum attenuation is
between 8.5 dB and 13.0 dB?

99
 b) According to the normal model, what proportion of the products has maximum
attenuation between 8.5 dB and 13.0 dB?
c) What proportion of the products has maximum attenuation greater than 15.1 dB?
Solution: Let X be the maximum attenuation of the next product, Then X is a normal variable
with μ=10.1 and σ=2.7.
X − 10.1
z=
2.7
a) Probability that the maximum attenuation of the next product is between 8.5 dB and
13.0 dB.
. . . .
=P(8.5≤X≤13.0)=P( ≤X≤ )=P(-0.59≤Z≤1.07)=P(Z≤1.07)-P(Z≤-0.59)
. .

=0.8577-0.2776=0.5801.
b) 0.5801 is the proportion of the product having maximum attenuation between 8.5 and
13.0 dB
c) Proportion of the products having maximum attenuation greater than 15.1 dB
=P(X>15.1)=P(Z>(15.1-10.1)/(2.7))=P(Z>1.85)=1-0.9678=0.0322
Example: 1 The actual amount of instant coffee that a filling machine puts into “4-ounce” jars
may be looked upon as a random variable having a normal distribution with σ=0.04 Ounce. If
only 2% of the jars are to contain less than 4 ounce, what should be the mean fill of these jars?
Solution: Let X be the amount in ounces of instant coffee that is put into jars. Then the
distribution of X is normal with mean μ and standard deviation 0.04 ounce.
We are given that P(X<μ)=0.02
Then P( < ) = 0.02
. .

P(Z< ) = 0.02 ⇨ = −2.05⇨μ=4.082 ounce.

. .

Session 30
Problems to be discussed by the faculty:
1. Given a Standard Normal distribution, find the area under the curve which lies
a) To the left of z=1.43;
b) to the right of z=-0.89
c) between z=-2.16 and z=-0.65
d) to the left of z=-1.39
e) to the right of z=1.96
f) between z=-0.48 and z= -1.74
2. Given a standard normal distribution find the value of k such that

100
a)𝑃(𝑧 < 𝑘) = 0.0427;
b)𝑃(𝑧 > 𝑘) = 0.2946;
c) P(-0.93<Z<k)=0.7235.
3. An electrical firm manufactures light bulbs that have a life, before burn out, i.e normally
distributed with mean of 800 hours and standard deviation of 40 hours. Find the
probability that a bulb (a) burns between 778 and 834 hours (b) burns out after 900
hours (c) burns out before 200 hours.
Home work:
1. A research scientist reports that mice will live on average of 40 months when their diets
are sharply restricted and then enriched with vitamins and proteins. Assuming that the
life times of such mice are normally distributed with a standard deviation of 6.3 months,
find the probability that a given mouse will live
a) more than 32 Months b) less than 28 months c) between 37 and 49 months.
Session 31
Problems to be discussed by the faculty:
1)In a photosnapic process, the developing time of prints may be looked upon as a random
variable having the normal distribution with mean 15.40 seconds and standard deviation 0.48
sec. Find the probabilities that the time it takes to develop one of the prints will be
a) atleast 16.00 sec
b) atmost 14.20 sec
c) anywhere from 15.00 to 15.80 sec
2) Suppose that the actual amount of instant coffee that a filling machine puts int6-ounce.
Jans is a random variable having a normal distribution with SD=0.05 ounce. If any 3% of
Jans are to contain less than 6 ounces of coffee,what must be the mean fill of these Jans?
3) A random variable has a normal distribution with SD=10. If the probability that the
random variable will take on a value less than 82.5% is 0.8212. What is the probability that it
will take on a value greater than 58.3?
Home Work
1) If 23% of all patients with high blood pressure have bad side effects from a certain kind of
medicine use the normal distribution to determine the probability that among 120 patients
with high blood pressure treated with this medicine more than 32 will have bad side effects?
2) The probability is 0.20 that a certain bank will refuse a loan application, use the normal
distribution to determine the probability that the bank will refuse atmost 40 of 225 loan
applications?
3)Suppose that during periods of transcendental meditation the reduction of a person’s
oxygen consumption with mean=37.6cc/min and SD =4.6CC/min. Find the probabilities that
during a period of transcendental meditation a person’s oxygen consumption will be reduced
by
a) atleast 44.5 CC.min
b) atmost 35.0CC/min
c) anywhere from 30.0 to 40.0CC/min
101
Session 32 INTRODUCTION TO MARKOV PROCESS
Stochastic process: A family of random variables which are functions of say, time are known
as stochastic process (or random process).
Eg 1. Consider the experiment of throwing an unbiased die. Suppose that X n is the outcome
of the nth throw, n  1. Then X n , n  1 is a family of random variables such that for distinct
values of n (=1, 2, 3…), one gets distinct random variable X n ; X n , n  1 constitutes a
stochastic process.
Eg 2. Suppose that X n is the number of sixes in the first n throws. For a distinct value of
n= 1,2,3……., we get a distinct binomial variable X n ; X n , n  1 which gives a family of
random variables is a stochastic process.
Eg 3. Suppose that X n is the maximum number shown in the first ‘n’ throws. Here X n , n  1
constitutes a stochastic process.
Eg 4. Consider the number of telephone calls received at a switch board. Suppose that X(t) the
random variable which represents the number of incoming calls in an interval (o, t) of duration
t units. The number of calls in one unit of time is X(1). The family { X (t ), t  T } constitutes a
stochastic process (T  [0,  )) .

State space: The set of possible values of a single random variable X n of a stochastic process
X n , n  1 is known as its state space. The state space is discrete if it contains a finite or
denumerable infinite of points; otherwise, it is continuous.
Eg 1. If X n is the total number of sixes appearing in the first n throws of a die, the set of
possible values of X n is the finite set of non negative integers 0,1,…n Hence the state space
of X n is discrete. We can write X n  y1  y 2  ...  y n , where yi is a discrete random variable
denoting the outcome of ith throw and yi=1 or 0 according as ith throw shows six or not.
Eg 2. Consider X n  z1  z 2  .....  z n , where zi is a continuous random variable assuming
values in [0, ). Here, the set of values of X n is the interval [0,  ), and so the state space
of X n is continuous.

Discrete time (parameter) stochastic process: In the above two examples we assume that the
parameter n of Xn is restricted to be non-negative integer, n=0,1,2,…. We consider the state of
system at discrete time points n = 0,1,2,…, only. Here the word time is used in a wide sense.
In the first example “time n” implies throw number n.
Continuous time (parameter) stochastic process: we can visualize a family of random
variables { Xt , t  T } (or { X (t ), t  T } ) such that the state of the system is characterized at
every instant over a finite or infinite interval. The system is then defined for a continuous
range of time and we say that we have a family of random variables in continuous time.
A stochastic process in continuous time may have either a discrete or continuous state
space.

102
Eg 1: Suppose that X(t) is the number of telephone calls at a switch board in an interval (o,
t). Here the state space of X(t) is discrete though X(t) is defined for a continuous range of time.
This is a continuous time stochastic process with discrete state space.
Eg 2. Suppose that X(t) represents the maximum temperature at a particular place in (0, t), then
set of possible values of X(t) is continuous. This is a continuous time stochastic process with
continuous state space.
Thus the stochastic processes can be classified into the following four types of processes:
(i) Discrete time; discrete state space
(ii) Discrete time; continuous state space
(iii) Continuous time; discrete state space
(iv) Continuous time; continuous state space

All the four types may be represented by { X (t ), t  T } . In case of discrete time, the
parameter generally used is n, i.e,. the family is represented by { X ( n ), n  0,1,2....} . In case
of continuous time both the symbols { X t , t  T } and { X (t ), t  T } ( where T is finite or
infinite interval) are used. The parameter t is usually interpreted as time, through it may
represent such characters as distance, length, thickness and so on.
Markov process: If { X (t ), t  T } is a stochastic process such that, for, t1  t 2  ...  t n  t

Pr a  X (t )  b / X (t1 )  x1 , X (t 2 )  x 2 ,... X (t n )  x n   Pra  X (t )  b / X (t n )  x n  the

process { X (t ), t  T } is a Markov process.

Markov chain: A discrete parameter Markov process is know as a Markov chain.

Definition: The stochastic process X n , n  0,1,2,3,... is called a Markov chain if, for
j , k , j1 ,  j n 1  N

Pr X n  k / X n 1  j , X n  2  j1 ,  , X 0  j n 1   PrX n  k / X n 1  j
,
whenever the first term is defined.
The outcomes are called the states of the Markov chain; If X n has the outcome j

(i.e. X n  j ) , the process is said to be at state j at the nth trail. To a pair of states (j, k) at
the two successive trails ( nth and (n+1)th trails) there is an associated conditional probability
Pjk .

Transition probability: Pjk is called the transition probability and represents the
probability of transition from state j at the nth trial to the state k at the (n+1)th trail.
Homogeneous Markov chain: If the transition probability Pjk is independent of n, the
Markov chain is said to be homogeneous. If it is dependent on n, the chain is said to be
non-homogeneous.

103
One step transition probability: The transition probability Pjk refer to the states (j, k) at
two successive trails (say nth and (n+1)th trails); the transition is one step transition
probability.
If we are concerned with the pair of states (j, k) at two non-successive trails, say, j at
the nth trail and k at the (n+m)th trail, the corresponding probability is then called m-step
 PrX n m  k / X n  j.
(m)
transition probability and is denoted by Pjk

Transition probability Matrix or Matrix of transition probabilities: The transition

probability Pjk satisfy Pjk >0,  Pjk  1 for all j. These probabilities may be written in the
matrix form

 p11 p12 p13 

p p 22 p 23 
 21
P   
 
   
 

This is called the transition probability Matrix of the Markov chain. P is a stochastic matrix.
Thus a transition matrix is a square matrix with non-negative elements and unit-row sums.
Example 1: Consider a communication system which transmits the two digits 0 and 1
through several stages. Let X n , n  1 be the digit leaving the nth stage of system and X0 be
the digit entering the first stage. At each stage there is a constant probability q that the digit
which enters will be transmitted unchanged, and probability p otherwise, p+q=1.
Here X n , n  0 is a homogeneous two-state Markov chain with single step transition
matrix
0 1

0 q p
P 
1p q 

Example 2. A particle performs a random walk with absorbing barriers say, as 0 and 4.
Whenever it is at any position r (0 < r < 4) , it moves to r+1 with probability p or to (r-1)
with probability q, p+q=1. But as soon as it reaches 0 or 4 it remains there itself. Let X n be
the position of the particle after n moves. The different states of Xn are the different
positions of the particle {Xn} is a Markov chain whose unit-step probabilities are given by

PrX n  r  1 / X n 1  r  p
PrX n  r  1 / X n 1  r  q, 0  r  4

and

PrX n  0 / X n 1  0  1
PrX n  4 / X n 1  4  1

104
The transition matrix is given by
States of X n

0 1 2 3 4

0 1 0 0 0 0
1 q 0 p 0 0 
States of X n 1 2 0 q 0 p 0
 
3 0 0 q 0 p
4  0 0 0 0 1 

Eg 3: General case of random walk: In the above example, as soon as the particle reaches
0 it remains there with probability ‘a’ and is reflected to state 1 with probability 1-a(0 < a
<1); if it reaches state 4 it remains there with probability b and reflected to state 3 with
probability 1-b (o < b<1) then {Xn} is a Markov chain with state space {0,1,2,3,4}.
The transition matrix is given by

0 1 2 3 4

0 a 1  a 0 0 0
1 q 0 p 0 0 
2 0 q 0 p 0
 
3 0 0 q 0 p
4  0 0 0 1 b b 

If a=1, then ‘0’ is an absorbing barrier and if a=0, then ‘0’ is a reflecting barrier similar is
a case with state 4.
Finite Markov chain: A Markov chain X n , n  0 with k states, where k is finite, is said
to be a finite Markov chain. The transition matrix P in this case is a square matrix with k
rows and k columns.
The number of states could however infinite. Where the possible values of X n form a
denumerable set, then the Markov chain is said to be denumerable infinite or denumerable
and the chain is said to have a countable state space.
Probability distribution:
The probability distribution of X 0 , X 1 , X 2 ,  , X n can be computed in terms of the
transition probability p jk and the initial distribution of X 0 .

105
 PrX 0  a, X 1  b,  , X n  2  i, X n 1  j , X n  k 
 PrX n  k / X n 1  j , X n  2  i,  , X 1  b, X 0  a
 PrX n  k / X n 1  j pr X n 1  j , X n  2  i,  , X 0  a
 PrX n  k / X n 1  jpr X n 1  j / X n  2  i PrX 1  b / X 0  aPrX 0  a
 p jk p ij  p ab PrX 0  a.

Eg: Let X n , n  0 be a Markov chain with three states 0,1,2 with transition matrix

0 1 2
0 3 1 
4 0
4
1  
1 1 1
2 4 2 4
 3 1 
0
 4 4 

1
and the initial distribution PrX 0  i  , i  0,1, 2.
3
3
We have PrX 1  1 / X 0  2 
4
1
PrX 2  2 / X 1  1 
4
PrX 2  2, X 1  1 / X 0  2
1 3 3
 PrX 2  2 / X 1  1PrX 1  1 / X 0  2  . 
4 4 16
3 1 1
PrX 2  2, X 1  1, X 0  2  PrX 2  2, X 1  1 / X 0  2PrX 0  2  .  .
16 3 16
PrX 3  1, X 2  2, X 1  1, X 0  2
3 1 3
PrX 3  1 / X 2  2Pr X 2  2, X 1  1, X 0  2  . 
4 16 64
Note: The transition probabilities with the initial distribution completely specifies a
Markov chain.
Higher order transition probabilities:

 PrX n  2  k / X n  j   PrX n  2  k , X n 1  r / X n  j
(2)
p jk
r

  p jr p rk
r

106
   p rk   p Jr
( m n) ( n) (m) (n) ( m)
p jk p jr p rk
r

Let p  ( p jk ) denote the transition matrix of the unit-step transition and P ( m )  Pjk ( m ) denote
the transition matrix of the m-step transitions. For m = 2, we have the matrix p ( 2)  p. p  p 2
similarly p ( m  n )  p ( m ) p  p. p ( m )

Classification of chains: The Markov chains are of two types (i) ergodic (ii) regular
An ergodic Markov chain has the property that it is possible to pass from one state to another
in a finite number of steps, regardless of present state.
A special type of ergodic Markov chain is the regular Markov chain.
A regular Markov chain is defined as a chain having a transition matrix P such that for some
power of P it has only non-zero positive probability values.
Note: Thus all regular chains must be ergodic chains.
Example: Consider a communication system which transmits the digits 0 and 1 through several
stages. At each stage the probability that the same digit will be received by the next stage, as
transmitted, is 0.75. What is the probability that a 0 that is entered at the first stage is received
as a 0 by the 5th stage?

Solution: Now we want to find the P004

0.75 0.25
The state transition matrix P is given by
P 
0.25 0.75
 0.625 0.375 
Hence P   0.375 0.625 
2

 

0.53125 0.46875
And
P 4
 P 2 2
P  0.46875 0.53125
 
Therefore the probability that a zero will be transmitted through four stages as a zero is
P004=0.53125.
Problems discussed by the faculty:
1. Classify the following states

1/ 4 3 / 4 
a) P 
1/ 2 1/ 2 

b)  1 0 
P 
1/ 4 3 / 4 

107
  0 0.5 0 0.5 
c)  0.25 0 0.75 0 
P 
 0 0.3 0 0.7 
 
 0.2 0 0.8 0 

2. A communication system transmits the two digits 0 and 1, each of them passing through
several stages. Suppose that the probability that the digit that enters remains unchanged
when it leaves, is p and that it changes is q  1  p . Suppose further that X 0 is the
digit which enters the first stage of the system and X n (n  1) is the digit leaving the nth
stage of the system. Show that { X n , n  1} forms a Markov chain. Find P, P2, P3 and
calculate PrX 2  0 / X 0  1 and PrX 3  1 / X 0  0

Home Work
1.Classify the states of the following Markov chains. If a state is periodic, determine its
period.

0.5 0.25 0.25 0

0 1 0
a) 0 0 1 b) 0 0 1 0
1/3 0 1/3 1/3
1 0 0
0 0 0 1

2. A communication system transmits the two digits 0 and 1, each of them passing through
several stages. Suppose that the probability that the digit that enters remains unchanged when
it leaves, is p and that it changes is q  1  p . Suppose further that X 0 is the digit which enters
the first stage of the system and X n (n  1) is the digit leaving the nth stage of the system. Show
that { X n , n  1} forms a Markov chain. Find P, P2, P3 and calculate PrX 2  0 / X 0  1 and
PrX 3  1 / X 0  0

3. The transition probability matrix of a Markov chain { X n , n  1,2 } having three states 1, 2
 0.1 0.5 0.4
and 3 is P  0.6 0.2 0.2 and the initial distribution is  0  (0.7, 0.2, 0.1) . Find
 0.3 0.4 0.3

(i) Pr X 2  3
(ii) Pr  X 3  2, X 2  3, X 1  3, X 0  2 .

Session 33 INTRODUCTION TO MARKOV PROCESS

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1. A market survey is made on two brands of breakfast food A and B. Every time a customer
purchases, he may buy the same brand or switch to another brand. The transition matrix is
given below:
To
A B

A 0.8 0.2
from
B 0.6 0.4

At present, it is estimated that 60% of the people buy brand A and 40% buy brand B.
Determine the market shares of brand A and brand B in the steady state.
2. An engineering professor acquires a new computer once every two years. The Professor
can choose from three models: M1, M2 and M3. If the present model is M1, the next
computer can be M2 with probability 0.2 or M3 with probability 0.15. If the present
model is M2, the probabilities of switching to M1 and M3 are 0.6 and 0.25 respectively.
And if the present model is M3, then the probabilities of purchasing M1 and M2 are 0.5
and 0.1 respectively. Represent the situation as a Markov chain.
Home Work
1. A house wife buys three kinds of cereals; A, B and C. She never buys the same cereal
on successive weeks. If she buys cereal A, then the next week she buys cereal B.
However, if she buys either B or C, then the next week she is three times as likely to
buy A as the other brand. Obtain the transition probability matrix and determine how
often she would buy each of the cereals in the long run.
2. A research analyzing brand switching between different airlines, operating on the
Delhi-Mumbai route by frequent fliers. On the basis of the data collected by her, the
researcher has developed the following transition probability matrix.
To airline

AA  0.9 .03 0.07

From airline BB 0.15 0.80 0.05
CC 0.20 0.30 0.50

It is found that currently the airlines AA, BB and CC have 20%, 50% and 30% of the
market respectively.
i) Obtain the market share for each airline in two moths time, and

Calculate the long run market share for each time.

Session 34
CO 4 Complex Variables
This course outcome covers the basic principles of differentiable complex-valued functions of
a single complex variable. Topics include the complex number system, Cauchy-Riemann

109
conditions, analytic functions and their properties, complex integration and line integrals,
Cauchy's theorem, Cauchy’s integral formula and simple applications.
The complex number system is merely a logical extension of the real number system. The
set of complex numbers includes the real numbers and still more. All complex numbers are
of the form x + iy where i =  1. In other words i2 = -1. If y = 0, then the complex number
x + iy becomes the real number x. Because there are two real numbers ( x and y ) associated
with each complex number, we are able to depict complex numbers using a plane, as opposed
to the real’s which are depicted on a line. Unlike the real number system, complex numbers
are not ordered. This means that it is not meaningful to say z1 < z2 in the complex number
system, even though such a thing is possible in the real’s.
It is possible to define addition and multiplication of complex numbers in the following
intuitive ways:
Addition: (x1+iy1) + (x2 + iy2) = (x1+x2) + i(y1+y2)
Multiplication: (x1+iy1)(x2+iy2) = (x1x2 – y1y2) + i(x1y2+x2y1)
The complex number 0 + i0 is the complex counterpart of zero in the real’s. It is the complex
additive identity. We will at times simply denote it as 0. The multiplicative identity is equal
to 1 + i0, which we will at times denote as 1.
A complex number can be written as z, so long as we understand that z = x + iy. It is
possible to discuss subtracting and dividing complex numbers. For example,
z1 – z2 = (x1+iy1) + (-x2 + i(-y2)) = (x1 - x2) + i(y1 - y2)

z1 x1 y
= ( x1  iy1 )( i 2 1 2 ) = 1
z2 x  y1
2
1
2
x1  y1

In addition to the basic operations of addition, subtraction, multiplication, and division, we

can also perform more complicated operations – such as taking the square root.

z1 = a + ib where (a+ib)(a+ib) = z1 = x1 + iy1 .

Example: Find 3  i4

Note that there are two solutions: 3  i 4 =2+i and 3  i 4 = -2-i

To check this we note that (2+i)(2+i) = 3 + i4 as is the case with –2-i.

It is not hard to show that there will be exactly two complex square roots for any (nonzero)
complex number.The complex conjugate of a complex number z = x + iy is denoted z and is
defined as (x – iy). The modulus of a complex number is defined as | z | zz .

Representation in the 2-Dimensional Plane

Each complex number can be written as z = x+iy. This means that we can associate an ordered
pair (x,y) with each and every complex number z = x+iy. Luckily, this gives us a graphical
representation of the complex number system. We can visualize the complex numbers. The 2-
dimensional plane that represents the complex numbers is sometimes called the Argand plane

110
but was first employed by Gauss. The horizontal axis represents the real numbers which is a
1-dimensional subspace of the plane. The vertical axis represents "pure" complex numbers; or
numbers which have no real part, x. A good question is whether or not the complex numbers
(field) is isomorphic to R2 under addition and multiplication.
The unit circle plays an important role in complex numbers since any (non-zero) complex
number can be written as a scalar multiple of its corresponding point on the unit circle. For
example, the point z = x + iy is a (x2 + y2) - multiple of
x y
i 2
x y
2 2
x  y2

which lies on the unit circle. Seen in this way, the unit circle can generate the entire set of
complex numbers through the appropriate multiplication of scalars to points on the unit circle.
This is analogous to a point on the real number line (e.g., 1 and –1 ) being able to generate any
other number on the real number line by the appropriate multiplication of a scalar. The unit
circle is not unique in this regard, though. Other types of geometric objects containing the
origin can do this as well. The value of the circle is that each point on the circle is equidistant
from the origin and this distance is equal to 1. Note that the real numbers 1 and –1 have
distance from zero equal to unity and can generate any real number through multiplication of
a scalar. It is interesting, in this regard, that the inverse of a complex number involves the
normalization of both the real and (negative) imaginary parts of the number. That is, the
denominator is the formula for a circle.
The unit circle separates the plane into two regions. The set of points that are strictly inside
the circle (called the open unit disk) and the set of points on and outside the unit circle. The
interior of the unit circle (i.e. the unit disk) is particularly important to the stability of certain
difference equations. It is similarly involved in determining whether a time series is covariance
stationary.
Functions of a Complex Variable
Functions of a complex variable provide us some powerful and widely useful tools in
mathematical analysis as well as in theoretical physics.
e.g.: • Some important physical quantities are complex variables
• Evaluating definite integrals.
• Obtaining asymptotic solutions of differentials equations.
• Integral transforms
• Many Physical quantities that were originally real become complex as simple theory
is made more general. The energy En  En0 + i (-1  the finite life time).
Note: A complex number z = (x,y) = x + iy
Where i = (-1)1/2. We will see that the ordering of two real numbers (x,y) is significant, i.e. in
general x + iy  y + ix
x is called the real part, labeled by Re z, y is called the imaginary part, labeled by Im z
For Cartesian components
z 1  z 2  x 1  x 2  i  y1  y 2 

111
 z 1 z 2  x 1 x 2  y 1 y 2   i x 1 y 2  x 2 y 1 

z=rei
r – the modulus or magnitude of z
 - the argument or phase of z

r  z  x 2
 y 
2 1/ 2

  tan  1  y / x 
The choice of polar representation or Cartesian representation is a matter of convenience.
Addition and subtraction of complex variables are easier in the Cartesian representation.
Multiplication, division, powers, roots are easier to handle in polar form,

z1 z 2  r1r2 e i 1  2 

z1 / z 2  r1 / r2 e i 1  2 

z n  r n e in
Using the vector analogy, we have the triangle inequalities

z1  z 2  z1  z 2  z1  z 2

Using the polar form,

z1 z 2  z1 z 2

arg(z1z2) = arg(z1) + arg(z2)

From z, complex functions f(z) may be constructed. They can be written
f(z) = u(x,y) + iv(x,y)
in which v and u are real. For example if f(z)=z2, we have

 
f  z   x 2  y 2  i 2 xy .

Complex Conjugation: replacing i by –i, which is denoted by (*),

z *  x  iy

We then have

zz *  x 2  y 2  r 2

hence

 
z  zz *
12

Note:

112
 z  re i or rei   2 n 

ln z  ln r  i or ln z  ln r  i  2n 

We can define complex valued functions of a complex variable. That

is, the domain of the function is the complex variable field and the range is also the complex
field. We can write this as
w = f(z)
where both w and z are complex numbers. Such functions have complex numbers as
parameters, as well. For example, we can write the following function
zo 1  i 2 ( x  2 y )  i 2( x  y )
w  f ( z)     u  iv
z x  iy x2  y2

Clearly, w is complex, as is z. The parameter zo is also complex, but it is a fixed complex

number. Note how that u and v have become real multivariate functions of x and y. That is,
x  2y 2( x  y )
u  u( x, y )  and v  v ( x, y ) 
x2  y2 x2  y2

As x and y run over all the values in R2, both u and v are determined, and hence z is
determined accordingly. This complex z then determines the value of w.
One of the most useful of all the complex functions is the exponential function. This function
has a straight forward relation to the trigonometric functions. We can understand this relation
by using a MacLaurin series for the ex function. To begin with

x x2
e 1 
x

1! 2!
Which, if we substitute iθ for x, we get

( i ) ( i ) 2
e i  1   
1! 2!

i 2  3i 4  5i 6
 1      
1! 2! 3! 4! 5! 6!

2 4 6 i  3i  5i  7i
 {1    }  {    }
2! 4! 6! 1! 3! 5! 7!
 cos( )  i sin( )

Note that the point ( cos(θ), sin(θ) ) is on the unit circle and as θ runs from 0 to 2π, the point
moves completely around the circle in a counterclockwise fashion. We can therefore write
any complex number as a scalar multiple of eiθ. Usually this is written as
z = x + iy = rcos(θ) + i rsin(θ) = reiθ
with θ = arctan(y/x) and r2 = x2 + y2.

113
The complex exponential and its relation to the trigonometric functions is of the greatest
imporance in of mathematics. It is incredibly useful and leads to some rather extraordinary
and unexpected results.
For example, it allows us to easily compute the following real number
1
z = ii   0.207
e

where i   1 . It also allows us to write out the logarithm of a negative number, which was
a great controversy during the time of Euler and Leibnitz. That is, we can write
z = ln(-1) = iπ
from which all other negative logarithms can be derived. 1 The logarithm of a complex
number can also be derived using this relation. Hence, we have
ln(z) = ln(x+iy) = ln(reiθ) = ln(r) + iθ
where θ is the angle formed by vectors (x, 0) and (0, y) and where r 2 = x2 + y2.
Polynomial equations, even simple ones, have solutions which are surprising to those who
look only for real solutions. For example, even the very simple equation
z4 + 1 = 0
has FOUR distinct roots (consider z2 = i and z2 = -i). In general, the Fundamental Theorem
of Algebra tells us that there will be exactly n complex roots (possibly repeated) which solve
an nth order polynomial equation. Once again, it is important to remember that the
coefficients on these polynomial equations can also be complex numbers, as well.
The familiar trigonometric functions of sin(x) and cos(x) can be defined for complex
numbers. This is done in the perfectly logical manner as follows:

e ix  e  ix e ix  e  ix
sin( z )  sin( x  iy )  and cos( z )  cos( x  iy )  ,
2i 2
1
where we remember that  i .
i
Limits and Derivatives of Complex Functions
Limits in the complex system are complicated by the fact that z = x + iy depends on (x, y)
and therefore one can approach zo = xo + iyo along infinitely many paths. For example,
1 1 1 1
z n  ( xo  )  i ( yo  ) and z n'  ( xo  )  i ( yo  )
n n n n
both limit to zo = xo + iyo, but do so along different paths. Obviously, other more
complicated paths are possible. This makes it a little more difficult to define a derivative,
which makes use of limits in its definition.

114
The derivative of w = f(z), if it exists, is defined by the unique limit
dw f ( z   )  f ( z)
 lim
dz   ( 0i 0 ) 
where   ( 0  i 0)  0 , the origin, along ANY path.

z
1. Is Ltexists? If so determine?
z
z 0

Example: w = f(z) = z2 is differentiable. To see this, assume (g(n), h(n)) are functions
parametrized such that they limit to the origin as n   . The ordered pair we have assumed
maps out any path to the origin and is perfectly general. It is not difficult to show that

{(x  g( n ))  i( y  h( n )}2  {x  iy}2

f ' ( z )  lim
n  g( n )  ih( n )

= lim [2{x  iy}  {g( n )  ih ( n )}] = 2z

n 

and thus, regardless of the path we take to the origin, the limit remains the same and thus the
derivative of f(z) is equal to f ’(z) = 2z. ■
Session 34
1. If z is a complex variable, then find the real and imaginary components of exp(z)
2. If z = 2+3i, then evaluate |exp(z)|
3. Given that e 2 i z  10 e i z  1  0 find e  y
4. If z=x+iy, then find the real and imaginary components of log (z)
5. Find the real and imaginary components of Sin(z)

HOMEWORK
1. Find , the real and imaginary components of exp(z2)
2. If w = f(z) = z2 + 3z, find u and v and then compute f(z) at z = 1 + 3i
3. If z=x+iy, then identify the real and imaginary components of Cos (z)
4. Evaluate the principal values of (1  i ) (1  i )
i
5. Solve for log ( z ) 
2
Session 35
The Cauchy-Riemann Equations and Complex Differentiation
Suppose that we consider f(z) = z2 and substitute into this z = x+iy. We can therefore write
this function again in the following way:
f(z) = z2 = F(x,y) = (x+iy)2 = (x2 – y2) + i2xy = u(x,y) + iv(x,y)
where u(x,y) = (x2 – y2) and where v(x,y) = 2xy. Now since z = x+iy, we know that
zz zz
x and y 
2 2i
x 1 y 1
from which it follows that  and  .
z 2 z 2 i

115
Now consider the complex derivative f ’(z).
F x F y 1 1
f ' ( z)    ( 2 x  2 yi )( )  ( 2 y  2 xi )( )
x z y z 2 2i

This of course reduces to f ' ( z )  2 z and the result agrees with the derivative computed in
the previous section using limits.
Now suppose that F(x,y) = u(x,y) + iv(x,y) is any differentiable complex function. What
must be true about the functions u and v ? This is the subject of the Cauchy - Riemann
equations.
First, suppose that z changes by x changing alone. Then, assume that z changes by y
changing alone. This would give us two expressions for the derivative of
f(z) = F(x,y).
The first (holding y constant) can be written as
F x u x v x 1 u v
f ' ( z ) | y cons tan t   i  { i }
x z x z x z 2 x x
while the second (holding x constant) can be written as
F y u y v y 1 u v
f ' ( z ) |x cons tan t   i  {i  }.
y z y z y z 2 y y

Now, the derivative of f(z) cannot depend on which way that z is changing (either by x
changing alone or alternatively by y changing alone ) and so the two expressions for f ' ( z )
must be equal if the derivative exists. This implies that
u v u v
 and that  
x y y x

These two equalities are known as the Cauchy-Riemann Equations.

Session 35
1. Verify the function e- x(cosy - i siny) is analytic or not?.
2. Does the function f ( z ) 2
 k where k is a constant satisfies the CR conditions
3. If u = sin (x) cosh (y) then verify, whether u satisfies the harmonic function or not.
4. If u = x2 - y2 - y then verify whether, u is harmonic or not. If harmonic then find its
conjugate. Then verify whether u and v satisfy CR conditions.
HOMEWORK
1. If u = ax3 +bxy, then under what conditions, 'a' and 'b' satisfy the harmonic
function property.
2. Verify the function f(z)=1/z2 is analytic or not?
3. If f(z)= sinh(z) then verify, whether the CR conditions are satisfied.
4. Verify that exp(z) is an entire function. Using CR equations.
5. Find the real and imaginary parts of an analytic function f(z)=exp(1/z 2).

116
Session 36 Construction of Analytic function by Milne-Thomson method:
Let f(z)=u+iv be analytic function
u v
f ( z )  i
x x
Use CR equations and then replace x by z and y by 0
After integration we obtain the required analytic function.
Session 36
Problems to be discussed by the faculty
1. Determine the function φ. If w=φ+iΨ represent the complex potential for an electric
field and Ψ= (x2-y2) + x/(x2+y2).
2. An electric field in the xy-plane is given by the potential function φ= 3x 2y-y2, find the
stream function ψ if exists.
3. The current in a conductor follows the relation u + i v. If the real part of the potential
function is given by u = e- ycosx, then find 'v'

Session 37
Problems to be discussed by the faculty
1. Find the velocity potential φ, in a two dimensional fluid flow. The stream
function Ψ= -y/x2+y2 is given.
2. Determine p such that the function f(z)= [loge(x2+y2)+itan-1(px/y)]/2 is analytic
and verify whether
imaginary component exists if u = [loge(x2+y2)]
3. Does the function e- xcos(y) satisfies the Laplace equation? If so find its
Harmonic conjugate
HOME WORK
1. The stream function in a potential flow is u  x 2  y 2 . Does 'v' exists ?If so
evaluate.
2. If u =(Coshx)Sin y can its counterpart 'v' can be evaluated. If so evaluate it.
3. Find u and v from w= exp(z2) and verify that the curves u(x,y)=c1 and v(x,y)=c2
cut orthogonally or not?.
Session 38
Introduction to Structure of Algebras, Semi groups, Monoids and Groups

Binary operation on a set: Let G is a non-empty set. Then G×G= {(a, b): a∈ 𝐺, 𝑏 ∈ 𝐺}. If f:
G× 𝐺 → 𝐺, then f is said to be a binary operation on the set G.
The image of the ordered pair (a, b) under f is denoted by a f b. Often we use symbols +,
×, . , 𝑜,∗ etc. to denote the binary operations on a set.
Thus ‘+’ will be a binary operation on G iff a+b∈ 𝐺, ∀𝑎, 𝑏 ∈ 𝐺 𝑎𝑛𝑑 𝑎 + 𝑏 𝑖𝑠 𝑢𝑛𝑖𝑞𝑢𝑒.
Similarly ‘*’ will be a binary operation on G iff a*b∈ 𝐺, ∀𝑎, 𝑏 ∈ 𝐺 𝑎𝑛𝑑 𝑎 ∗ 𝑏 𝑖𝑠 𝑢𝑛𝑖𝑞𝑢𝑒.
Note: A binary operation on a set G is sometimes also called as binary composition in the set
G. If ‘*’ is a binary composition in G then a*b∈ 𝐺, ∀𝑎, 𝑏 ∈ 𝐺. Therefore G is closed with
respect to the composition denoted by ‘*’.

117
Example: Addition is a binary operation on set N of natural numbers. The sum of two natural
numbers is also a natural number. Therefore N is closed with respect to addition. i.e., a+b∈
𝑁, ∀𝑎, 𝑏 ∈ 𝑁.
Subtraction is not a binary operation on N, we have 4-7=-3∉N, whereas 4∈ 𝑁, 7 ∈ 𝑁.
Thus N is not closed with respect to subtraction.
But subtraction is a binary operation on the set of integers I. We have a-b∈ 𝐼, ∀𝑎, 𝑏 ∈ 𝐼.
Algebraic Structure: Definition: A non-empty set G equipped with one or more binary
operations is called an algebraic structure.
Suppose * is a binary operation on G, then (G, *) is an algebraic structure.
(N, +), (I, +), (I, -), (R, +, .) are all algebraic structures. Obviously addition and multiplication
are both binary operations on the set R of real numbers. Therefore (R, +, .) is an algebraic
structure equipped with two operations.
Semi Group: Definition: An algebraic structure (G, *) is called a semi group if the binary
operation * is associative in G. i.e. if (a*b)*c=a*(b*c) ∀𝑎, 𝑏, 𝑐 ∈ 𝐺.
Example: The set N of all natural numbers is a semi group with respect to the operation of
addition of natural numbers. Obviously addition is an associative operation on N.
Similarly the algebraic structures (N, .), (I, +), (R, +) are also semi groups.
Group: Definition: Let G be a non-empty set equipped with a binary operation denoted by ‘.’
Then this algebraic structure (G, .) is a group, if the binary operation ‘.’ Satisfies the following
properties:
1. Closure Property: i.e., ab∈ 𝐺, ∀𝑎, 𝑏 ∈ 𝐺.
2. Associativity: i.e., (ab)c=a(bc) ∀𝑎, 𝑏, 𝑐 ∈ 𝐺
3. Existence of Identity: There exists an element e∈ 𝐺 ∋ 𝑒𝑎 = 𝑎 = 𝑎𝑒 ∀𝑎 ∈ 𝐺. The
element e is called the identity element.
4. Existence of Inverse: Each element of G possesses inverse. In other words a∈
𝐺 ⟹there exists an element b∈ 𝐺 such that ba=e=ab. The element b is called the
inverse of a and we write b=𝑎 . Thus 𝑎 is an element of G such that 𝑎 𝑎 = 𝑎𝑎 =
𝑒.

Abelian or Commutative group: A group G is said to be abelian or commutative if in addition

to the above four properties the following property is also satisfied:
5. Commutativity: i.e., a b = b a ∀𝑎, 𝑏 ∈ 𝐺.

Finite and infinite groups:

If in a group G, the set G consisting of a finite number of distinct elements then the group is
called a finite group, otherwise an infinite group. The number of elements in a finite group is
called the order of the group. An infinite group is said to be of infinite order.
We shall denote the order of the group G denoted by the symbol o(G)

Problems to be discussed by the faculty

1. Is the set I of all integers ……..-3, -2, -1, 0, 1, 2, 3,……is a group with respect to the
operation of addition of integers?
2. Determine whether the set N of all natural numbers 1, 2, 3, 4, ……..is a group with
respect to addition.

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 3. Is the set 𝑄 of all non-zero rational numbers forms a group under the operation of
multiplication of rational numbers?
4. Is the set G = {1, -1} is an abelian group w.r.t. multiplication of real numbers?

Homework problems:
1. Show that the set of positive rational numbers forms an abelian group under the
composition defined by a*b=a b/2.
2. Is the set I of integers …..-3, -2, -1, 1, 2, 3,…..is a group (i) with respect to
subtraction (ii) with respect to multiplication ?
3. Is the set of even natural numbers is a group (i) under addition (ii) under
multiplication?
4. Let 𝑄 be the set of all positive rational numbers and * a binary operation on 𝑄
defined by a*b= . Determine the identity element in 𝑄 and determine the inverse
of a.
5. Is the set G = {1, -1, i, -i} is an abelian group w.r.t. multiplication of Complex
numbers?

Session 39

Homorphisms :-
The concept of homomorphism for algebraic systems are define as follows
Definition:- Let (X,.) and (Y,*) be two algebraic systems of the same type in the sence that
both . and * are binary (n-ary) operations. A mapping g: X  Y is called a homomorphism,
or simply morphism, from (X,.) to (y,*) if for any x 1,x2  X , g(x1.x2)=g(x1)*g(x2).
If such a function g exits, then it is customary to call (Y,*) a homomorphic image of (X,.).
Now we apply this concept to semigroups, monoids and groups. Homomorphisms of
semigroups and monoids have useful applications in the economical design of sequential
machines and in formal languages.
Definition:- Let (S,  ) and (T,  ) be any semigroups. A mapping g: S  T such that for any
two
elements a,b  S, g(a  b) = g(a)  g(b) is called a semigroup homomorphism.
Definition:- Let (M,  ,em) and (N,  ,en) be any two monoids. A mapping g:M  N such
that for any two elements a,b  M (i) g(a  b) = g(a)  g(b) (ii) g(em) = en. is called a monoid
homomorphism. (Here em and en are identity elements in monoids M and N respectively.).
Definition:- Let (G,  ) and (H,  ) be two groups . A mapping g:G  His called a group
homomoephism from (G,  ) to (H,  ) if for any a,b  G g(a  b) = g(a)  g(b).
Epimorphisms, and Monomorphisms :--

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Definition :- Let G and H be groups. Then, a function g : G → H is a epihomomorphism iff
g is surjective homomorphism .
Definition :- Let G and H be groups. Then, a function g : G → H is a monohomomorphism
iff g is injective homomorphism .
Isomorphism:-
Definition:- Let G and H be groups. Then, a function g : G → H is a isomorphism iff g is
injective and surjective homomorphism .
Endomorphism:-
Definition :- Let G be a group. An homomorphism g : G → G is called an endomorphism.
We write Endo(G) = {g : G → G | g is an homomorphism.}.
Automorphism:-
Definition :- Let G be a group. An isomorphism g : G → G is called an automorphism. We
write Auto(G) = {g : G → G | g is an automorphism.}.
Examples:-
1.) Let (G,+) be the additive group of integers and (G*.) be the multiplicative group whose
elements are in the form 2m for m   . Consider the mapping g:G  G* by g(n) = 2n for n  
.Then g is a homomorphism from G to G*.

To verify,
For any m, n  G, by closer property m+n  G.
Also 2m and 2n are elements of G* hence by closer property 2m.2n= 2m+n  G.

Then by definition of mapping, g(m+n) = 2m+n = 2m.2n= g(m).g(n).

Hence g is homomorphism from G into G*.
  
2.) Let (G,+) be the additive group of integers, G' = Zm = { 0, 1,....... m 1 }be the group of
resedue classes modulo m w.r.t addition of residue classes.

Define the mapping g:G  G' by g(a) = a for a  G.
Further Every a  G has a unique has a unique image in G'. Similarly every element of G' is
the image of some element in G. Hence g is onto mapping.
  
Also for a,b  G , g(a+b) = a  b = a  b = g(a)+g(b).
Therefore g is a homomorphism.
Hence g is called epimorphism.

3.) Let (  , +) be semigroup of natural numbers and ( S ,  ) be the semigroup on S = {e,0,1}

With the operation  given by Table .
 e 0 1
e e 0 1
0 0 0 0
1 1 0 1

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A mapping g:   S , given by g(0)=1 and g(j)=0 for j  0 is a semigroup homomorphism.

Note:- Although both (  , +) and ( S ,  ) are monoids with identities o and e respectively, the
mapping g is not a monoid homomorphism, because g(0)  e.

4.) Let (G,+) be a group of real numbers under addition and (H,.) is a group of positive real
numbers under multiplication. Let f:G  H be a mapping such that f(x)=ex for x  G.
Show
That f is Isomorphisim.
Solution:- If x is a real number, ex>0 and hence ex  H.

Let a,b  G Then, ea,eb  H. Now f(a) = f(b)  ea = eb  a = b.

Therefore f is one-one mapping

Let c  H, i.e., c is a positive real number and logc is a real number, so logc  G.

Therefore f(logc) = elogc = c. Thus there exits logc  G such that f(logc) = c.

Hence f is on-to mapping.

Further , Let a,b  G by closer property a+b  G , then f(a+b) = ea+b= ea.eb = f(a).f(b).

Therefore f is homomorphism.
From these f is one—one, onto, homomorphism mapping. Hence f is Isomorphism.

Section 39
Problems to be discussed by the faculty

1.) Let g: S  T be an isomorphism of semigroups ( S , ) and (T,  ) , if e is zero

element of S , Is g(e) zero element of ( T ,  ).
2.) Let G be the multiplicative group of all 2x2 non-singular matrices whose elements
are real numbers and G’ be the multiplicative group of non-zero real numbers. Show
that the mapping  :G  G’ by  (A) = |A|for A  G is homomorphism.
3.) Verify the mapping f:G  G’ such that f(x+iy) = x where G is a group of complex
numbers under addition and G’ is a group of real numbers is isomorphism.
Session 40

Normal subgroups and congruence Relations:

Normal subgroup: Definition: A sub group H of a group G is said to be a normal sub group
of G if for every x∈ 𝐺 and for every h∈ 𝐻, xh𝑥 ∈ 𝐻.
Note: From this definition we can conclude that H is a normal subgroup of G iff xH𝑥 ⊆ 𝐻
∀𝑥 ∈ 𝐺.
Congruence relations:

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Relation of ‘congruence modulo m’ in the set of integers:
Let m be any fixed positive integer. Then an integer a is said to be congruent to another integer
b modulo m if m|a-b i.e., if a-b is divisible by m. symbolically we write a≡ 𝑏 (𝑚𝑜𝑑 𝑚). It will
be read as “a is congruent to b modulo m”.
Example 1: 15≡3(mod 6) because 15-3=12 which is divisible by 6.
Example 2: 23≡-3(mod 2) because 23-(-3)=26 which is divisible by 2.
Problems to be discussed by the faculty
1. Verify that the ‘congruence modulo m’ is an equivalence relation in the set of integers?
2. Is 17≡-3(mod 5)?

Homework problems:
1. Is 13≡3(mod 5)?
2. Is 17≡3(mod 5)?

Session 41
RINGS

Definition : An algebraic system S ,, is called a ring if the binary operations  and  on
S satisfy the following three properties :

1. S , is an abelian group.
2. S , is a semigroup.
3. The operation  is distributive over  ; that is, for any a, b, c  S ,
a  b  c   a  b  a  c and b  c   a  b  a  c  a
Familiar examples of rings are the sets of integers, real numbers, rational numbers,
even numbers, and complex numbers under the operations of addition and multiplication.
Because of these examples, it is customary to refer to the operation  as addition and  as
multiplication in a ring S ,, although these operations may not necessarily ean additions
and multiplications. In keeping with this convention, we shall also refer to the identity of
S , as the additive identity and denote it by 0. Similarly, if S , is a monoid, then the
identity with respect to  will be called the multiplicative identity and will be denoted by 1.
The additive inverse of an element a  S will be denoted by  a , while the multiplicative
inverse, if it exists, will be denoted by a 1 . It must be emphasized at this point that the use of
this terminology does not mean that the operations  and  on the ring S ,, have all the
properties that  and  have in the system of real numbers.

Depending upon the structure of S , , various special cases of rings will be defined.

If S , is commutative, then the ring S ,, is called a commutative ring. Similarly,

if S , is a monoid, then S ,, is called a ring with identity.

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Example 1 : Prove that the set of even integers is a ring w.r.t usual addition and
multiplication of integers.
Solution: Let R  the set of even integers.

Then R  2 x / x  Z .

a , b , c  R  a  2 m, b  2 n , c  2 p where m, n, p  Z .

a  b  2( m  n ) is even integer.

 addition   of integers in a binary operation in R .

Additive identity is 0.
Further, every element of a  R has additive inverse as  a .
In general integers addition is commutative.

Thus R,  is a commutative group.

a  b  2m2n   2l where l  2mn  Z .

 Multiplication  of integers in a binary operation in R .

a  b  c  2m  2n   2 p  8mnp .
a  b  c   2m  2n  2 p   8mnp .

 a  b  c  a  b  c  .
 Multiplication  is associative in R .

a  b  c   2m2n  2 p 

 2 m  2n  2m  2 p

 ab  ac
Similarly, b  c   a  b  a  c  a .

 Distributive laws hold in R .

Hence R,, is a ring.

Example 2: Prove that the set of all 2  2 matrices over the complex numbers is a ring
with unity under addition and multiplication of matrices.
𝑎 𝑏
Sol. Let 𝑎 , 𝑏, 𝑐, 𝑑 𝑎𝑟𝑒 𝑐𝑜𝑚𝑝𝑙𝑒𝑥 𝑛𝑢𝑚𝑏𝑒𝑟𝑠 and A, B , C  R .
𝑐 𝑑
Since A, B are 2  2 matrices, by the definition the sum of A and B i.e., A  B is also
a 2  2 matrix whose elements are complex numbers.
 addition of matrices is a binary operation in R .

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 For A, B  R we have A  B  B  A , since the complex numbers commute under
addition.
As addition of matrices is associative,
 A  B   C  A  B  C A, B, C  R.
 0 0
There exists 0     R such that A  0  A for each A  R .
 0 0
a b  a  b 
For A     R we have  A     R such that A   A  0 .
c d  c  d
 R,  is an abelian group.
Since A, B are 2  2 matrices, by the definition the product of A and B i.e., A  B is
also a 2  2 matrix whose elements are complex numbers.
 multiplication of matrices is a binary operation in R . As multiplication of matrices
is associative,
 AB C  ABC A, B, C  R .
Since multiplication of matrices is distributive over addition
AB  C   AB  AC and B  C A  BA  CAA, B, C  R .
Hence R,, is a ring.

Example 3 : The algebraic system Z n , n , n consisting of equivalence classes generated

by the relation congruence modulo n for some fixed integer n over the set of integers is a
ring. The operations  n and  n were correspond to the operations of addition and
multiplication over integers.

Observe that for n  6 , Z 6 ={ 0, 1 , 2 , 3 , 4 , 5 } then Z 6 , 6 , 6 is a ring.

Example 4 : Let S be a set and  S  its power set. On  S  we define operations  and 

as follows:

A  B  x  S / x  A  x  B    x  A  B 

A B  A B

for all A, B   S  . It is easy to verify that  S ,, is a ring called the ring of subsets of S .

Using the definition of a ring one can prove the following results for a ring R,, in
which 0 is the additive identity and  a denotes the additive inverse of an element a  R .
1) a  0  0  a  0
2) a   b   a  b 
3)  a   b  a  b
4)  a    b  a  b
5) a  b  c   a  b  a  c
6) a  b  c  a  c  b  c

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 In 5) we write b  c in place of b   c  and a  b  a  c for a  b   a  c  .

Definition : A subset R  S where S ,, is a ring is called a subring if R,, is itself a

ring with the operations  and  restricted to R .
Example : The ring of even integers is a subring of the ring of integers.
In fact, if R  S , if we determine that R is closed with respect to addition and that for
any a  R, a  R , and finally if R is closed with respect to the operation  , then R is a
subring of S . All other properties of a ring are satisfied by R . It may happen that R may not
have some additional properties which S may possess. For example, note that the ring of
even integers is not a ring with identity, while the ring of integers is a ring with identity.
Problems to be discussed by the faculty
1. If R= a, b, c, d ,, is a ring whose operations are given by the following table :

+ a b c d . a b c d

a a b c d a a a a a
b b c d a b a c a c
c c d a b c a a a a
d d a b c d a c a a

Is it a commutative ring? Find the additive inverse of each of its elements.

 
2. If z 2 = {a+b 2 /a,b  z }is a ring with respect to addition and multiplication of
numbers.( z is set of all integers). Determine the additive identity and additive
inverse
3. If I ,, is a ring , where the operations  and  are defined,
=
as a  b  a  b  1 and a  b  a  b  ab . For every a, b   .Verify the
ring I is commutative ring with unity are not.
4. Prove that if a, b  R where R,, is a ring, then
a  b 2  a 2  a  b  b  a  b 2 where a 2  a  a .

*******

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